Measure, Integral and Probability

Marek Capiński and Ekkehard Kopp Measure, Integral and Probability Springer-Verlag Berlin Heidelberg New York London ...

Author: Marek Capinski & Ekkhard Kopp

47 downloads 1609 Views 1MB Size Report

This content was uploaded by our users and we assume good faith they have the permission to share this book. If you own the copyright to this book and it is wrongfully on our website, we offer a simple DMCA procedure to remove your content from our site. Start by pressing the button below!
Report copyright / DMCA form

DOWNLOAD PDF

Marek Capiński and Ekkehard Kopp

Measure, Integral and Probability

Springer-Verlag Berlin Heidelberg New York London Paris Tokyo Hong Kong Barcelona Budapest

To our children; grandchildren: Piotr, Maciej, Jan, Anna; Lukasz Anna, Emily

Preface

The central concepts in this book are Lebesgue measure and the Lebesgue integral. Their role as standard fare in UK undergraduate mathematics courses is not wholly secure; yet they provide the principal model for the development of the abstract measure spaces which underpin modern probability theory, while the Lebesgue function spaces remain the main source of examples on which to test the methods of functional analysis and its many applications, such as Fourier analysis and the theory of partial differential equations. It follows that not only budding analysts have need of a clear understanding of the construction and properties of measures and integrals, but also that those who wish to contribute seriously to the applications of analytical methods in a wide variety of areas of mathematics, physics, electronics, engineering and, most recently, finance, need to study the underlying theory with some care. We have found remarkably few texts in the current literature which aim explicitly to provide for these needs, at a level accessible to current undergraduates. There are many good books on modern probability theory, and increasingly they recognize the need for a strong grounding in the tools we develop in this book, but all too often the treatment is either too advanced for an undergraduate audience or else somewhat perfunctory. We hope therefore that the current text will not be regarded as one which fills a much-needed gap in the literature! One fundamental decision in developing a treatment of integration is whether to begin with measures or integrals, i.e. whether to start with sets or with functions. Functional analysts have tended to favour the latter approach, while the former is clearly necessary for the development of probability. We have decided to side with the probabilists in this argument, and to use the (reasonably) systematic development of basic concepts and results in probability theory as the principal field of application – the order of topics and the vii

viii

Preface

terminology we use reflect this choice, and each chapter concludes with further development of the relevant probabilistic concepts. At times this approach may seem less ‘efficient’ than the alternative, but we have opted for direct proofs and explicit constructions, sometimes at the cost of elegance. We hope that it will increase understanding. The treatment of measure and integration is as self-contained as we could make it within the space and time constraints: some sections may seem too pedestrian for final-year undergraduates, but experience in testing much of the material over a number of years at Hull University teaches us that familiarity and confidence with basic concepts in analysis can frequently seem somewhat shaky among these audiences. Hence the preliminaries include a review of Riemann integration, as well as a reminder of some fundamental concepts of elementary real analysis. While probability theory is chosen here as the principal area of application of measure and integral, this is not a text on elementary probability, of which many can be found in the literature. Though this is not an advanced text, it is intended to be studied (not skimmed lightly) and it has been designed to be useful for directed self-study as well as for a lecture course. Thus a significant proportion of results, labelled ‘Proposition’, are not proved immediately, but left for the reader to attempt before proceeding further (often with a hint on how to begin), and there is a generous helping of Exercises. To aid self-study, proofs of the Propositions are given at the end of each chapter, and outline solutions of the Exercises are given at the end of the book. Thus few mysteries should remain for the diligent. After an introductory chapter, motivating and preparing for the principal definitions of measure and integral, Chapter 2 provides a detailed construction of Lebesgue measure and its properties, and proceeds to abstract the axioms appropriate for probability spaces. This sets a pattern for the remaining chapters, where the concept of independence is pursued in ever more general contexts, as a distinguishing feature of probability theory. Chapter 3 develops the integral for non-negative measurable functions, and introduces random variables and their induced probability distributions, while Chapter 4 develops the main limit theorems for the Lebesgue integral and compares this with Riemann integration. The applications in probability lead to a discussion of expectations, with a focus on densities and the role of characteristic functions. In Chapter 5 the motivation is more functional-analytic: the focus is on the Lebesgue function spaces, including a discussion of the special role of the space L2 of square-integrable functions. Chapter 6 sees a return to measure theory, with the detailed development of product measure and Fubini’s theorem, now leading to the role of joint distributions and conditioning in probability. Finally,

Preface

ix

following a discussion of the principal modes of convergence for sequences of integrable functions, Chapter 7 adopts an unashamedly probabilistic bias, with a treatment of the principal limit theorems, culminating in the Lindeberg–Feller version of the Central Limit Theorem. The treatment is by no means exhaustive, as this is a textbook, not a treatise. Nonetheless the range of topics is probably slightly too extensive for a one-semester course at third-year level: the first five chapters might provide a useful course for such students, with the last two left for self-study or as part of a reading course for students wishing to continue in probability theory. Alternatively, students with a stronger preparation in analysis might use the first two chapters as background material and complete the remainder of the book in a one-semester course. May 1998

Marek Capi´ nski Ekkehard Kopp

Preface to the Second Edition

After five years and six printings it seems only fair to our readers that we should respond to their comments and also correct errors and imperfections to which we have been alerted in addition to those we have discovered ourselves in reviewing the text. This second edition also introduces additional material which earlier constraints of time and space had precluded, and which has, in our view, become more essential as the make-up of our potential readership has become clearer. We hope that we manage to do this in a spirit which preserves the essential features of the text, namely providing the material rigorously and in a form suitable for directed self-study. Thus the focus remains on accessibility, explicitness and emphasis on concrete examples, in a style that seeks to encourage readers to become directly involved with the material and challenges them to prove many of the results themselves (knowing that solutions are also given in the text!). Apart from further examples and exercises, the new material presented here is of two contrasting kinds. The new Chapter 7 adds a discussion of the comparison of general measures, with the Radon-Nikodym Theorem as its focus. The proof given here, while not new, is in our view more constructive and elementary than the usual ones, and we utilise the result consistently to examine the structure of Lebesgue-Stieltjes measures on the line and to deduce the Hahn-Jordan decomposition of signed measures. The common origin of the concepts of variation and absolute continuity of functions and measures is clarified. The main probabilistic application is to conditional expectations, for which an alternative construction via orthogonal projections is also provided in Chapter 5. This is applied in turn in Chapter 7 to derive elementary properties of martingales in discrete time. The other addition occurs at the end of each chapter (with the exception of Chapters 1 and 5). Since it is clear that a significant proportion of our current xi

xii


readership is amongst students of the burgeoning field of mathematical finance, each relevant chapter ends with a brief discussion of ideas from that subject. In these sections we depart from our aim of keeping the book self-contained, since we can hardly develop this whole discipline afresh. Thus we neither define nor explain the origin of the finance concepts we address, but instead seek to locate them mathematically within the conceptual framework of measure and probability. This leads to conclusions with a mathematical precision that sometimes eludes authors writing from a finance perspective. To avoid misunderstanding we repeat that the purpose of this book remains the development of the ideas of measure and integral, especially with a view to their applications in probability and (briefly) in finance. This is therefore neither a textbook in probability theory nor in mathematical finance. Both of these disciplines have a large specialist literature of their own, and our comments on these areas of application are intended to assist the student in understanding the mathematical framework which underpins them. We are grateful to those of our readers and to colleagues who have pointed out many of the errors, both typographical and conceptual, of the first edition. The errors that inevitably remain are our sole responsibility. To facilitate their speedy correction a webpage has been created for the notification of errors, inaccuracies and queries, at http://www.springer.co.uk/MIP and we encourage our readers to use it mercilessly. Our thanks also go to Stephanie Harding and Karen Borthwick at Springer Verlag, London, for their continuing care and helpfulness in producing this edition. October 2003

Marek Capi´ nski Ekkehard Kopp

Contents

1.

Motivation and preliminaries . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1 1.1 Notation and basic set theory . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2 1.1.1 Sets and functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2 1.1.2 Countable and uncountable sets in R . . . . . . . . . . . . . . . . . 4 1.1.3 Topological properties of sets in R . . . . . . . . . . . . . . . . . . . . 5 1.2 The Riemann integral: scope and limitations . . . . . . . . . . . . . . . . . 7 1.3 Choosing numbers at random . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 12

2.

Measure . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.1 Null sets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.2 Outer measure . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.3 Lebesgue measurable sets and Lebesgue measure . . . . . . . . . . . . . 2.4 Basic properties of Lebesgue measure . . . . . . . . . . . . . . . . . . . . . . . 2.5 Borel sets . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.6 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.6.1 Probability space . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2.6.2 Events: conditioning and independence . . . . . . . . . . . . . . . . 2.6.3 Applications to mathematical finance . . . . . . . . . . . . . . . . . 2.7 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

15 15 20 26 35 40 45 46 46 49 51

3.

Measurable functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.1 The extended real line . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.2 Lebesgue-measurable functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.3 Examples . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.4 Properties . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.5 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

55 55 55 59 60 66

xiii

xiv

Contents

3.5.1 Random variables . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.5.2 Sigma fields generated by random variables . . . . . . . . . . . . 3.5.3 Probability distributions . . . . . . . . . . . . . . . . . . . . . . . . . . . . 3.5.4 Independence of random variables . . . . . . . . . . . . . . . . . . . . 3.5.5 Applications to mathematical finance . . . . . . . . . . . . . . . . . 3.6 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .

66 67 68 70 71 73

4.

Integral . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 75 4.1 Definition of the integral . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 75 4.2 Monotone Convergence Theorems . . . . . . . . . . . . . . . . . . . . . . . . . . 82 4.3 Integrable functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 86 4.4 The Dominated Convergence Theorem . . . . . . . . . . . . . . . . . . . . . . 92 4.5 Relation to the Riemann integral . . . . . . . . . . . . . . . . . . . . . . . . . . . 97 4.6 Approximation of measurable functions . . . . . . . . . . . . . . . . . . . . . 102 4.7 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 105 4.7.1 Integration with respect to probability distributions . . . . 105 4.7.2 Absolutely continuous measures: examples of densities . . 106 4.7.3 Expectation of a random variable . . . . . . . . . . . . . . . . . . . . . 114 4.7.4 Characteristic function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 115 4.7.5 Applications to mathematical finance . . . . . . . . . . . . . . . . . 117 4.8 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 119

5.

Spaces of integrable functions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 125 5.1 The space L1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 126 5.2 The Hilbert space L2 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 131 5.2.1 Properties of the L2 -norm . . . . . . . . . . . . . . . . . . . . . . . . . . . 132 5.2.2 Inner product spaces . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 135 5.2.3 Orthogonality and projections . . . . . . . . . . . . . . . . . . . . . . . 137 5.3 The Lp spaces: completeness . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 140 5.4 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 146 5.4.1 Moments . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 146 5.4.2 Independence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 150 5.4.3 Conditional Expectation (first construction) . . . . . . . . . . . 153 5.5 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 155

6.

Product measures . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 159 6.1 Multi-dimensional Lebesgue measure . . . . . . . . . . . . . . . . . . . . . . . . 159 6.2 Product σ-fields . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 160 6.3 Construction of the product measure . . . . . . . . . . . . . . . . . . . . . . . . 162 6.4 Fubini’s Theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 169 6.5 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 173 6.5.1 Joint distributions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 173

Contents

xv

6.5.2 Independence again . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 175 6.5.3 Conditional probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 178 6.5.4 Characteristic functions determine distributions . . . . . . . . 180 6.5.5 Application to mathematical finance . . . . . . . . . . . . . . . . . . 182 6.6 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 185 7.

The Radon–Nikodym Theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 187 7.1 Densities and Conditioning . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 187 7.2 The Radon–Nikodym Theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 188 7.3 Lebesgue–Stieltjes measures . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 198 7.3.1 Construction of Lebesgue–Stieltjes measures . . . . . . . . . . . 199 7.3.2 Absolute continuity of functions . . . . . . . . . . . . . . . . . . . . . . 204 7.3.3 Functions of bounded variation . . . . . . . . . . . . . . . . . . . . . . . 206 7.3.4 Signed measures . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 210 7.4 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 218 7.4.1 Conditional expectation relative to a σ-field . . . . . . . . . . . 218 7.4.2 Martingales . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 221 7.4.3 Applications to mathematical finance . . . . . . . . . . . . . . . . . 231 7.5 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 234

8.

Limit theorems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 241 8.1 Modes of convergence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 241 8.2 Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 243 8.2.1 Convergence in probability . . . . . . . . . . . . . . . . . . . . . . . . . . 245 8.2.2 Weak law of large numbers . . . . . . . . . . . . . . . . . . . . . . . . . . 249 8.2.3 The Borel–Cantelli Lemmas . . . . . . . . . . . . . . . . . . . . . . . . . 255 8.2.4 Strong law of large numbers . . . . . . . . . . . . . . . . . . . . . . . . . 260 8.2.5 Weak convergence . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 268 8.2.6 Central Limit Theorem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 273 8.2.7 Applications to mathematical finance . . . . . . . . . . . . . . . . . 280 8.3 Proofs of propositions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 283

9.

Solutions to exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 287

10. Appendix . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 301 References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 305 Bibliography . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 305 Index . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 307

1 Motivation and preliminaries

Life is an uncertain business. We can seldom be sure that our plans will work out as we intend, and are thus conditioned from an early age to think in terms of the likelihood that certain events will occur, and which are ‘more likely’ than others. Turning this vague description into a probability model amounts to the construction of a rational framework for thinking about uncertainty. The framework ought to be a general one, which enables us equally to handle situations where we have to sift a great deal of prior information, and those where we have little to go on. Some degree of judgement is needed in all cases; but we seek an orderly theoretical framework and methodology which enables us to formulate general laws in quantitative terms. This leads us to mathematical models for probability, that is to say, idealized abstractions of empirical practice, which nonetheless have to satisfy the criteria of wide applicability, accuracy and simplicity. In this book our concern will be with the construction and use of generally applicable probability models in which we can also consider infinite sample spaces and infinite sequences of trials: that such are needed is easily seen when one tries to make sense of apparently simple concepts such as ‘drawing a number at random from the interval [0, 1]’ and in trying to understand the limit behaviour of a sequence of identical trials. Just as elementary probabilities are computed by finding the comparative sizes of sets of outcomes, we will find that the fundamental problem to be solved is that of measuring the ‘size’ of a set with infinitely many elements. At least for sets on the real line, the ideas of basic real analysis provide us with a convincing answer, and this contains all the ideas needed for the abstract axiomatic framework on which to base the theory of probability. For 1

2


this reason the development of the concept of measure, and Lebesgue measureon R in particular, has pride of place in this book.

1.1 Notation and basic set theory In measure theory we deal typically with families of subsets of some arbitrary given set and consider functions which assign real numbers to sets belonging to these families. Thus we need to review some basic set notation and operations on sets, as well as discussing the distinction between countably and uncountably infinite sets, with particular reference to subsets of the real line R. We shall also need notions from analysis such as limits of sequences, series, and open sets. Readers are assumed to be largely familiar with this material and may thus skip lightly over this section, which is included to introduce notation and make the text reasonably self-contained and hence useful for self-study. The discussion remains quite informal, without reference to foundational issues, and the reader is referred to basic texts on analysis for most of the proofs. Here we mention just two recent introductory textbooks: [8] and [11].

1.1.1 Sets and functions In our operations with sets we shall always deal with collections of subsets of some universal set Ω; the nature of this set will be clear from the context – frequently Ω will be the set R of real numbers or a subset of it. We leave the concept of ‘set’ as undefined and given, and concern ourselves only with set membership and operations. The empty set is denoted by Ø; it has no members. Sets are generally denoted by capital letters. Set membership is denoted by ∈, so x ∈ A means that the element x is a member of the set A. Set inclusion, A ⊂ B, means that every member of A is a member of B. This includes the case when A and B are equal; if the inclusion is strict, i.e. A ⊂ B and B contains elements which are not in A (written x ∈ / A) this will be stated separately. The notation {x ∈ A : P (x)} is used to denote the set of elements of A with property P . The set of all subsets of A (its power set) is denoted by P(A). We define the intersection A ∩ B = {x : x ∈ A and x ∈ B} and union A ∪ B = {x : x ∈ A or x ∈ B}. The complement Ac of A consists of the elements of Ω which are not members of A; we also write Ac = Ω \ A, and, more generally, we have the difference B \ A = {x ∈ B : x ∈ / A} = B ∩ Ac and the symmetric difference A∆B = (A \ B) ∪ (B \ A). Note that A∆B = Ø if

1. Motivation and preliminaries

3

and only if A = B. Intersection (resp. union) gives expression to the logical connective ‘and’ (resp. ‘or’) and, via the logical symbols ∃ (there exists) and ∀ (for all), they have extensions to arbitrary collections; indexed by some set Λ these are given by \ Aα = {x : x ∈ Aα for all α ∈ Λ} = {x : ∀α ∈ Λ, x ∈ Aα } α∈Λ

[

α∈Λ

Aα = {x : x ∈ Aα for some α ∈ Λ} = {x : ∃α ∈ Λ, x ∈ Aα }.

These are linked by de Morgan’s laws: [ \ ( A α )c = Acα ; α

α

(

\ α

A α )c =

[

Acα .

α

If A ∩ B = Ø then A and B are disjoint. A family of sets (Aα )α∈Λ is pairwise disjoint if Aα ∩ Aβ = Ø whenever α 6= β (α, β ∈ Λ). The Cartesian product A × B of sets A and B is the set of ordered pairs A × B = {(a, b) : a ∈ A, b ∈ B}. As already indicated, we use N, Z, Q, R for the basic number systems of natural numbers, integers, rationals and reals respectively. Intervals in R are denoted via each endpoint, with a square bracket indicating its inclusion, an open bracket exclusion, e.g. [a, b) = {x ∈ R : a ≤ x < b}. We use ∞ and −∞ to describe unbounded intervals, e.g. (−∞, b) = {x ∈ R : x < b}, [0, ∞) = {x ∈ R : x ≥ 0} = R+ . R2 = R×R denotes the plane, more generally, Rn is the n-fold Cartesian product of R with itself, i.e. the set of all n-tuples (x1 , . . . , xn ) composed of real numbers. Products of intervals, called rectangles, are denoted similarly. Formally, a function f : A → B is a subset of A × B in which each first coordinate determines the second: if (a, b), (a, c) ∈ f then b = c. Its domain Df = {a ∈ A : ∃b ∈ B, (a, b) ∈ f }, and range Rf = {b ∈ B : ∃a ∈ A, (a, b) ∈ f } describe its scope. Informally, f associates elements of B with those of A, such that each a ∈ A has at most one image b ∈ B. We write this as b = f (a). The set X ⊂ A has image f (X) = {b ∈ B : b = f (a) for some a ∈ X} and the inverse image of a set Y ⊂ B is f −1 (Y ) = {a ∈ A : f (a) ∈ Y }. The composition f2 ◦ f1 of f1 : A → B and f2 : B → C is the function h : A → C defined by h(a) = f2 (f1 (a)). When A = B = C, x 7→ (f1 ◦ f2 )(x) = f1 (f2 (x)) and x 7→ (f2 ◦ f1 )(x) = f2 (f1 (x)) both define functions from A to A. In general, these will not be the same: for example, let f1 (x) = sin x, f2 (x) = x2 , then x 7→ sin(x2 ) and x 7→ (sin x)2 are not equal. The function g extends f if Df ⊂ Dg and g = f on Df ; alternatively we say that f restricts g to Df . These concepts will be used frequently for real-valued set functions, where the domains are collections of sets and the range is a subset of R.

4


The algebra of real functions is defined pointwise, i.e. the sum f + g and product f · g are given by (f + g)(x) = f (x) + g(x), (f · g)(x) = f (x) · g(x). The indicator function 1A of the set A is the function 1 for x ∈ A 1A (x) = 0 for x ∈ / A. Note that 1A∩B = 1A · 1B , 1A∪B = 1A + 1B − 1A 1B , and 1Ac = 1 − 1A . We need one more concept from basic set theory, which should be familiar: For any set E, an equivalence relation on E is a relation (i.e. a subset R of E × E, where we write x ∼ y to indicate that (x, y) ∈ R) with the following properties: 1. reflexive: for all x ∈ E, x ∼ x,

2. symmetric: x ∼ y implies y ∼ x,

3. transitive: x ∼ y and y ∼ z implies x ∼ z. An equivalence relation ∼ on E partitions E into disjoint equivalence classes: given x ∈ E, write [x] = {z : z ∼ x} for the equivalence class of x, i.e. the set of all elements of E that are equivalent to x. Thus x ∈ [x], hence S E = x∈E [x]. This is a disjoint union: if [x] ∩ [y] 6= Ø, then there is z ∈ E with x ∼ z and z ∼ y, hence x ∼ y, so that [x] = [y]. We shall denote the set of all equivalence classes so obtained by E/ ∼.

1.1.2 Countable and uncountable sets in R We say that a set A is countable if there is a one–one correspondence between A and a subset of N, i.e. a function f : A → N that takes distinct points to distinct points. Informally, A is finite if this correspondence can be set up using only an initial segment {1, 2, ..., N } of N (for some N ∈ N), while we call A countably infinite or denumerable if all of N is used. It is not difficult to see that countable unions of countable sets are countable; in particular, the set Q of rationals is countable. Cantor showed that the set R cannot be placed in one–one correspondence with (a subset of) N; thus it is an example of an uncountable set. Cantor’s proof assumes that we can write each real number uniquely as a decimal (always choosing the non-terminating version). We can also restrict ourselves (why?) to showing that the interval [0,1] is uncountable. If this set were countable, then we could write its elements as a sequence (xn )n≥1 , and since each xn has a unique decimal expansion of the form xn = 0.an1 an2 an3 ....ann ...


5

for digits aij chosen from the set {0, 1, 2..., 9}, we could therefore write down the array x1 = 0.a11 a12 a13 . . . x2 = 0.a21 a22 a23 . . . x3 = 0.a31 a32 a33 . . . ... Now write down y = 0.b1 b2 b3 . . ., where the digits bn are chosen to differ from ann . Such a decimal expansion defines a number y ∈ [0, 1] that differs from each of the xn (since its expansion differs from that of xn in the nth place). Hence our sequence does not exhaust [0,1], and the contradiction shows that [0,1] cannot be countable. Since the union of two countable sets must be countable, and since Q is countable, it follows that R\Q is uncountable, i.e. there are far ‘more’ irrationals than rationals! One way of making this seem more digestible is to consider the problem of choosing numbers at random from an interval in R. Recall that rational numbers are precisely those real numbers whose decimal expansion recurs (we include ‘terminates’ under ‘recurs’). Now imagine choosing a real number from [0,1] at random: think of the set R as a pond containing all real numbers, and imagine you are ‘fishing’ in this pond, pulling out one number at a time. How likely is it that the first number will be rational, i.e. how likely are we to find a number whose expansion recurs? It would be like rolling a ten-sided die infinitely many times and expecting, after a finite number of throws, to say with certainty that all subsequent throws will give the same digit. This does not seem at all likely, and we should therefore not be too surprised to find that countable sets (including Q) will be among those we can ‘neglect’ when measuring sets on the real line in the ‘unbiased’ or uniform way in which we have used the term ‘random’ so far. Possibly more surprising, however, will be the discovery that even some uncountable sets can be ‘negligible’ from the point of view adopted here.

1.1.3 Topological properties of sets in R Recall the definition of an open set O ⊂ R :

6


Definition 1.1 A subset O of the real line R is open if it is a union of open intervals, i.e. for intervals (Iα )α∈Λ , where Λ is some index set (countable or not) [ O= Iα . α∈Λ

A set is closed if its complement is open. Open sets in Rn (n > 1) can be defined as unions of n-fold products of intervals. This definition seems more general than it actually is, since, on R, countable unions will always suffice – though the freedom to work with general unions will be convenient later on. If Λ is an index set and Iα is an open interval for each α ∈ Λ, then there exists a countable collection (Iαk )k≥1 of these intervals whose union equals ∪α∈Λ Iα . What is more, the sequence of intervals can be chosen to be pairwise disjoint. It is easy to see that a finite intersection of open sets is open; however, a countable intersection of open sets need not be open: let On = (− n1 , 1) for n ≥ 1, then E = ∩∞ n=1 On = [0, 1) is not open. Note that R, unlike Rn or more general spaces, has a linear order, i.e. given x, y ∈ R we can decide whether x ≤ y or y ≤ x. Thus u is an upper bound for a set A ⊂ R if a ≤ u for all a ∈ A, and a lower bound is defined similarly. The supremum (or least upper bound) is then the minimum of all upper bounds and written sup A. The infimum (or greatest lower bound) inf A is defined as the maximum of all lower bounds. The completeness property of R can be expressed by the statement that every set which is bounded above has a supremum. A real function f is said to be continuous if f −1 (O) is open for each open set O. Every continuous real function defined on a closed bounded set attains its bounds on such a set, i.e. has a minimum and maximum value there. For example, if f : [a, b] → R is continuous, M = sup{f (x) : x ∈ [a, b]} = f (xmax ), m = inf{f (x) : x ∈ [a, b]} = f (xmin ) for some points xmax , xmin ∈ [a, b]. The Intermediate Value Theorem says that a continuous function takes all intermediate values between the extreme ones, i.e. for each y ∈ [m, M ] there is a θ ∈ [a, b] such that y = f (θ). Specializing to real sequences (xn ), we can further define the upper limit lim supn xn as inf{ sup xm : n ∈ N} m≥n

and the lower limit lim inf n xn as sup{ inf xm : n ∈ N}. m≥n


7

The sequence xn convergesif and only if these quantities coincide and their common value is then its limit. Recall that a sequence (xn ) converges and the real number x is its limit, written x = limx→∞ xn , if for every ε > 0 there is P an N ∈ N such that |xn − x| < ε whenever n ≥ N. A series n≥1 an converges Pm if the sequence xm = n=1 an of its partial sums converges, and its limit is P then the sum ∞ n=1 an of the series.

1.2 The Riemann integral: scope and limitations In this section we give a brief review of the Riemann integral, which forms part of the staple diet in introductory analysis courses, and consider some of the reasons why it does not suffice for more advanced applications. Let f : [a, b] → R be a bounded real function, where a, b, with a < b, are real numbers. A partition of [a, b] is a finite set P = {a0 , a1 , a2 , . . . , an } with a = a0 < a1 < a2 < . . . < an = b. The partition P gives rise to the upper and lower Riemann sums U (P, f ) =

n X

Mi ∆ai ,

L(P, f ) =

i=1

where ∆ai = ai − ai−1 ,

Mi =

n X

mi ∆ai

i=1

sup

f (x)

ai−1 ≤x≤ai

and mi =

inf

ai−1 ≤x≤ai

f (x)

for each i ≤ n. (Note that Mi and mi are well-defined real numbers since f is bounded on each interval [ai−1 , ai ].) In order to define the Riemann integral of f , one first shows that for any given partition P , L(P, f ) ≤ U (P, f ), and next that for any refinement, i.e. a partition P 0 ⊃ P , we must have L(P, f ) ≤ L(P 0 , f ) and U (P 0 , f ) ≤ U (P, f ). Finally, since for any two partitions P1 and P2 , their union P1 ∪ P2 is a refinement of both, we see that L(P, f ) ≤ U (Q, f ) for any partitions P, Q. The set {L(P, f ) : P is a partition of [a, R bb]} is thus bounded above in R, and we call its supremum the lower integral a f of f on [a, b]. Similarly, the infimum of the Rb set of upper sums is the upper integral a f . The function f is now said to be

8


Riemann-integrable on [a, b] if these two numbers R b coincide, and their common value is the Riemann integral of f , denoted by a f or, more commonly, Z

b

f (x) dx. a

This definition does not provide a convenient criterion for checking the integrability of particular functions; however, the following formulation provides a useful criterion for integrability – see [8] for a proof.

Theorem 1.1 (Riemann’s Criterion) f : [a, b] → R is Riemann-integrable if and only if for every ε > 0 there exists a partition Pε such that U (Pε , f ) − L(Pε , f ) < ε.

Example 1.1 R1 √ We calculate 0 f (x) dx when f (x) = x: our immediate problem is that square roots are hard to find except for perfect squares. Therefore we take partition points which are perfect squares, even though this means that the interval lengths of the different intervals do not stay the same (there is nothing to say that they should do, even if it often simplifies the calculations). In fact, take the sequence of partitions 2 i 1 Pn = {0, ( )2 , ( )2 , . . . , ( )2 , . . . , 1} n n n and consider the upper and lower sums, using the fact that f is increasing: U (Pn , f ) =

L(Pn , f ) = Hence

n n X i i−1 2 1 X 2 i ) }= 3 (2i − i) ( ){( )2 − ( n n n n i=1 i=1

n n X i−1 i i−1 2 1 X 2 ( ){( )2 − ( ) }= 3 (2i − 3i + 1). n n n n i=1 i=1

U (Pn , f ) − L(Pn , f ) =

n 1 X 1 1 (2i − 1) = 3 {n(n + 1) − n} = . 3 n i=1 n n

By choosing n large enough, we can make this difference less than any given ε > 0, hence f is integrable. The integral must be 32 , since both U (Pn , f ) and L(Pn , f ) converge to this value, as is easily seen.


9

Riemann’s criterion still does not give us a precise picture of the class of Riemann-integrable functions. However, it is easy to show (see [8]) that any bounded monotone function belongs to this class, and only a little more difficult to see that any continuous function f : [a, b] → R (which is of course automatically bounded) will be Riemann-integrable. This provides quite sufficient information for many practical purposes, and the tedium of calculations such as that given above can be avoided by proving

Theorem 1.2 (Fundamental Theorem of Calculus) If f : [a, b] → R is continuous and the function F : [a, b] → R has derivative f (i.e. F 0 = f on (a, b)) then Z b F (b) − F (a) = f (x) dx. a

This result therefore links the Riemann integral with differentiation, and displays F as a primitive (also called ‘anti-derivative’) of f : Z x F (x) = f (x) dx −a

up to a constant, thus justifying the elementary techniques of integration that form part of any Calculus course. We can relax the continuity requirement. A trivial step is to assume f bounded and continuous on [a, b] except at finitely many points. Then f is Riemann integrable. To see this split the interval into pieces on which f is continuous. Then f is integrable on each and hence one can derive integrability of f on the whole interval. As an example consider a function f equal to zero for all x ∈ [0, 1] except a1 , . . . , an where it equals 1. It is integrable with integral over [0, 1] equal to 0. Taking this further, however, will require the power of the Lebesgue theory: in Theorem 4.23 we show that f is Riemann-integrable if and only if it is continuous at ‘almost all’ points of [a, b]. This result is by no means trivial, as you will discover if you try to prove directly that the following function f , due to Dirichlet, is Riemann-integrable over [0, 1]: 1 if x = m n n ∈Q f (x) = 0 if x ∈ / Q. In fact, it is not difficult, see [8], to show that f is continuous at each irrational and discontinuous at every rational point, hence (as we will see) is continuous at ‘almost all’ points of [0, 1].

10


Since the purpose of this book is to present Lebesgue’s theory of integration, we should discuss why we need a new theory of integration at all: what, if anything, is wrong with the simple Riemann integral described above? First, scope: it doesn’t deal with all the kinds of functions that we hope to handle. The results that are most easily proved rely on continuous functions on bounded intervals; in order to handle integrals over unbounded intervals, e.g. Z ∞ 2 e−x dx −∞

or the integral of an unbounded function: Z 1 1 √ dx, x 0 we have to resort to ‘improper’ Riemann integrals, defined by a limit process: e.g. considering the integrals Z n Z 1 2 1 √ dx, e−x dx, x −n ε and letting n → ∞ or ε → 0 respectively. This isn’t all that serious a flaw. Second, dependence on intervals: we have no easy way of integrating over more general sets, or of integrating functions whose values are distributed ‘awkwardly’ over sets that differ greatly from intervals. For example, consider the upper and lower sums for the indicator function 1Q of Q over [0, 1]; however we partition [0, 1], each subinterval must contain both rational and irrational points; thus each upper sum is 1 and each lower sum 0. Hence we cannot calculate the Riemann integral of f over the interval [0,1]; it is simply ‘too discontinuous’. (You may easily convince yourself that f is discontinuous at all points of [0, 1].) Third, lack of completeness: rather more importantly from the point of view of applications, the Riemann integral doesn’t interact well with taking the limit of a sequence of functions. One may expect results of the following form: if a sequence fn of Riemann-integrable functions converges (in some appropriate Rb Rb sense) to f , then a fn dx → a f dx. We give two counterexamples showing what difficulties can arise if the functions (fn ) converge to f pointwise, i.e. fn (x) → f (x) for all x. 1. The limit need not be Riemann integrable, and so the convergence question does not even make sense. Here we may take f = 1Q , fn = 1An where


11

An = {q1 , . . . , qn }, and the sequence (qn ), n ≥ 1 is an enumeration of the rationals, so that (fn ) is even monotone increasing. 2. The limit is Riemann integrable, but the convergence of Riemann integrals does not hold. Let f = 0, consider [a, b] = [0, 1], and put  1 2 if 0 ≤ x < 2n  4n x 1 fn (x) = 4n − 4n2 x if 2n ≤ x < n1  1 0 if n ≤ x ≤ 1.

Figure 1.1 Graph of fn . This is a continuous function with integral 1. On the other hand, the sequence fn (x) converges to f = 0 since for all x, fn (x) = 0 for n sufficiently large (such that n1 < x). See Figure 1.1. To avoid problems of this kind, we can introduce the idea of uniform convergence: a sequence (fn ) in C[0, 1] converges uniformly to f if the sequence an = sup{|fn (x) − f (x)| : 0 ≤ x ≤ 1} converges to 0. In this case one can easily prove the convergence of the Riemann integrals: Z 1 Z 1 fn (x) dx → f (x) dx. 0

0

However, the ‘distance’ sup{|f (x) − g(x)| : 0 ≤ x ≤ 1} has nothing to do with integration as such and the uniform convergence is too restrictive for many R1 applications. A more natural concept of ‘distance’, given by 0 |f (x) − g(x)| dx, leads to another problem. Defining  if 0 ≤ x ≤ 12  0 1 gn (x) = n(x − 2 ) if 12 < x < 12 + n1  1 otherwise R1 it can be shown that 0 |gn (x) − gm (x)| dx → 0 as m, n → ∞; in Figure 1.2 the shaded area vanishes. (We say that (fn ) is a Cauchy sequence in this distance.)

12


Yet there is no continuous function f to which this sequence converges since the pointwise limit is f (x) = 1 for x > 12 and 0 otherwise, so that f = 1( 21 ,1] . So the space C([0, 1]) of all continuous functions f : [0, 1] → R is too small from this point of view.

1 g

m

g n

1 2

1 1 + 2 m

1 1 + 2 n

Figure 1.2 Graphs of gn , gm This is rather similar to the situation which leads one to work with R rather than just with the set of rationals Q (there are Cauchy sequences √ without limits in Q, for example a sequence of rational approximations of 2). Recalling the crucial importance of completeness in the case of R, we naturally look for a theory of integration which does not have this shortcoming. In the process we shall find that our new theory, which will include the Riemann integral as a special case, also solves the other problems listed.

1.3 Choosing numbers at random Before we start to develop the theory of Lebesgue measure to make sense of the ‘length’ of a general subset of R, let us pause to consider some practical motivation. The simplicity of elementary probability with finite sample spaces vanishes rapidly when we have an infinite number of outcomes, such as when we ‘pick a number between 0 and 1 at random’. We face making sense of the ‘probability’ that a given x ∈ [0, 1] is chosen. A similar, slightly more general question, is the following: what is the probability that the number we pick is rational? First a prior question: what do we mean by saying that we pick the number x at random? ‘Random’ plausibly means that in each such trial, each real number is ‘equally likely’ to be picked, so that we impose the uniform probability distribution on [0, 1]. But the ‘number’ of possible choices is infinite. Hence the event Ax that a fixed x is chosen ought to have zero probability. On the other


13

hand, since some number between 0 and 1 is chosen, and it is not impossible that it could be our x. Thus a set Ax 6= Ø can have P (Ax ) = 0. Our way of ‘measuring’ probabilities need not, therefore, be able to distinguish completely between sets – we could not really expect this in general if we want to handle infinite sets. We can go slightly further: the probability that any one of a finite set of reals A = {x1 , x2 , . . . , xn } is selected should also be 0, since it seems natural Pn that this probability P (A) should equal i=1 P ({xi }). We can extend this to claim the finite additivity property of the probability function A 7→ P (A), i.e. Sn Pn that if A1 , A2 , ..., An are disjoint sets, then P ( i=1 Ai ) = i=1 P (Ai ). This claim looks very plausible, and we shall see that it becomes an essential feature of any sensible basis for a calculus of probabilities. Less obvious is the claim that, under the uniform distribution, any countably infinite set, such as Q, must also carry probability 0 – yet that is exactly what an analysis of the ‘area under the graph’ of the function 1Q suggests. We can reinterpret this as a result of a ‘continuity property’ of the mapping A 7→ P (A) when we let n → ∞ in the above: if the sequence (Ai ) of subsets of R is disjoint then we would like to have P(

∞ [

i=1

Ai ) = lim

n→∞

n X i=1

P (Ai ) =

∞ X

P (Ai ).

i=1

We shall see in Chapter 2 that this condition is indeed satisfied by Lebesgue measure on the real line R, and it will be used as the defining property of abstract measures on arbitrary sets. There is much more to probability than is developed in this book: for example, we do not discuss finite sample spaces and the elegant combinatorial ideas that characterize a good introduction to probability, such as [6] and [9]. Our focus throughout remains on the essential role played by Lebesgue measure in the description of probabilistic phenomena based on infinite sample spaces. This leads us to leave to one side many of the interesting examples and applications which can be found in these texts, and provide, instead, a consistent development of the theoretical underpinnings of random variables with densities.

2 Measure

2.1 Null sets The idea of a ‘negligible’ set relates to one of the limitations of the Riemann integral, as we saw in the previous chapter. Since the function f = 1Q takes a non-zero value only on Q, and equals 1 there, the ‘area under its graph’ (if such makes sense) must be very closely linked to the ‘length’ of the set Q. This is why it turns out that we cannot integrate f in the Riemann sense: the sets Q and R \ Q are so different from intervals that it is not clear how we should measure their ‘lengths’ and it is clear that the ‘integral’ of f over [0, 1] should equal the ‘length’ of the set of rationals in [0, 1]. So how should we define this concept for more general sets? The obvious way of defining the ‘length’ of a set is to start with intervals nonetheless. Suppose that I is a bounded interval of any kind, i.e. I = [a, b], I = [a, b), I = (a, b] or I = (a, b). We simply define the length of I as l(I) = b−a in each case. As a particular case we have l({a}) = l([a, a]) = 0. It is then natural to say that a one-element set is ‘null’. Before we extend this idea to more general sets, first consider the length of a finite set. A finite set is not an interval but since a single point has length 0, adding finitely many such lengths together should still give 0. The underlying concept here is that if we decompose a set into a finite number of disjoint intervals, we compute the length of this set by adding the lengths of the pieces. As we have seen, in general it may not be always possible actually to decom15

16


pose a set into intervals. Therefore, we consider systems of intervals that cover a given set. We shall generalize the above idea by allowing a countable number of covering intervals. Thus we arrive at the following more general definition of sets of ‘zero length’:

Definition 2.1 A null set A ⊆ R is a set that may be covered by a sequence of intervals of arbitrarily small total length, i.e. given any ε > 0 we can find a sequence {In : n ≥ 1} of intervals such that A⊆ and

∞ X

∞ [

In

n=1

l(In ) < ε.

n=1

(We also say simply that ‘A is null ’.)

Exercise 2.1 Show that we get an equivalent notion if in the above definition we replace the word ‘intervals’ by any of these: ‘open intervals’, ‘closed intervals’, ‘the intervals of the form (a, b], ‘the intervals of the form [a, b)’. Note that the intervals do not need to be disjoint. It follows at once from the definition that the empty set is null. Next, any one-element set {x} is a null set. For, let ε > 0 and take I1 = (x − 4ε , x + 4ε ), In = [0, 0] for n ≥ 2. (Why take In = [0, 0] for n ≥ 2? Well, why not! We could equally have taken In = (0, 0) = Ø, of course!) Now ∞ X

n=1

l(In ) = l(I1 ) =

ε < ε. 2

More generally, any countable set A = {x1 , x2 , ...} is null. The simplest way to show this is to take In = [xn , xn ], for all n. However, as a gentle introduction to the next theorem we will cover A by open intervals. This way it is more fun.

2. Measure

17

For, let ε > 0 and cover A with the following sequence of intervals 1 1 ε· 1 2 2 1 1 ε ε I2 = (x2 − 16 , x2 + 16 ) l(I2 ) = ε · 2 2 2 1 1 ε ε I3 = (x3 − 32 , x3 + 32 ) l(I3 ) = ε · 3 2 2 ... ... 1 1 In = (xn − 2·2ε n , xn + 2·2ε n ) l(In ) = ε · n 2 2 I1 = (x1 − 8ε , x1 + 8ε )

Since

P∞

1 n=1 2n

= 1,

∞ X

l(In ) =

n=1

l(I1 ) =

ε 0. Our goal is to cover the set N by countably many intervals with total length less than ε. The proof goes in three steps, each being a little bit tricky. Step 1. We carefully cover each Nn by intervals. ‘Carefully’ means that the lengths have to be small. ‘Small’ means that we are going to add them up later to end up with a small number (and ‘small’ here means less than ε). Since N1 is null, there exist intervals Ik1 , k ≥ 1, such that ∞ X

k=1

l(Ik1 )
0, there is z ∈ ZA , z < ε. But a member of ZA is the total length of some covering of A. That is, there is a covering (In ) of A with total length less than ε, so A is null. This combines our general outer measure with the special case of ‘zero measure’. Note that m∗ (Ø) = 0, m∗ ({x}) = 0 for any x ∈ R, and m∗ (Q) = 0 (and in fact, for any countable X, m∗ (X) = 0). Next we observe that m∗ is monotone: the bigger the set, the greater its outer measure.

Proposition 2.3 If A ⊂ B then m∗ (A) ≤ m∗ (B). Hint Show that ZB ⊂ ZA and use the definition of inf.

22


The second step is to relate outer measure to the length of an interval. This innocent result contains the crux of the theory, since it shows that the formal definition of m∗ , which is applicable to all subsets of R, coincides with the intuitive idea for intervals, where our thought processes began. We must therefore expect the proof to contain some hidden depths, and we have to tackle these in stages: the hard work lies in showing that the length of the interval cannot be greater than its outer measure: for this we need to appeal to the famous Heine–Borel theorem, which states that every closed, bounded subset B of R is compact: given any collection of open sets Oα covering B (i.e. S B ⊂ α Oα ), there is a finite subcollection (Oαi )i≤n which still covers B, i.e. Sn B ⊂ i=1 Oαi (for a proof see [1]).

Theorem 2.4 The outer measure of an interval equals its length.

Proof If I is unbounded, then it is clear that it cannot be covered by a system of intervals with finite total length. This shows that m∗ (I) = ∞ and so m∗ (I) = l(I) = ∞. So we restrict ourselves to bounded intervals. Step 1. m∗ (I) ≤ l(I). We claim that l(I) ∈ ZI . Take the following sequence of intervals: I1 = I, In = [0, 0] for n ≥ 2. This sequence covers the set I, and the total length is equal to the length of I hence l(I) ∈ ZI . This is sufficient since the infimum of ZI cannot exceed any of its elements. Step 2. l(I) ≤ m∗ (I). (1) I = [a, b]. We shall show that for any ε > 0 l([a, b]) ≤ m∗ ([a, b]) + ε.

(2.1)

This is sufficient since we may obtain the required inequality passing to the limit, ε → 0. (Note that if x, y ∈ R and y > x then there is an ε > 0 with y > x + ε, e.g. ε = 12 (y − x).) So we take an arbitrary ε > 0. By the definition of outer measure we can find a sequence of intervals In covering [a, b] such that ∞ X

ε l(In ) ≤ m∗ ([a, b]) + . 2 n=1

(2.2)

2. Measure

23

We shall slightly increase each of the intervals to an open one. Let the endpoints of In be an , bn , and we take ε ε Jn = an − n+2 , bn + n+2 . 2 2 It is clear that

l(In ) = l(Jn ) − so that

ε 2n+1

∞ X

l(In ) =

∞ X

l(Jn ) ≤ m∗ ([a, b]) + ε.

n=1

We insert this in (2.2) and we have

n=1

∞ X

ε l(Jn ) − . 2 n=1

(2.3)

The new sequence of intervals of course covers [a, b] so by the Heine–Borel theorem we can choose a finite number of Jn to cover [a, b] (the set [a, b] is compact in R). We can add some intervals to this finite family to form an initial segment of the sequence (Jn ) – just for simplicity of notation. So for some finite index m we have [a, b] ⊆

m [

Jn .

(2.4)

n=1

Let Jn = (cn , dn ). Put c = min{c1 , . . . , cm }, d = max{d1 , . . . , dm }. The covering (2.4) means that c < a and b < d hence l([a, b]) < d − c. Next, the number d − c is certainly smaller than the total length of Jn , n = 1, . . . , m (some overlapping takes place) and l([a, b]) < d − c
0. ε ε l( (a, b) ) = l([a + , b − ]) + ε 2 2 ε ε ≤ m∗ ([a + , b − ]) + ε (by (1)) 2 2 ≤ m∗ ( (a, b) ) + ε (by Proposition 2.3).

24


(3) I = [a, b) or I = (a, b]. l(I) = l((a, b)) ≤ m∗ ((a, b)) ∗

≤ m (I)

(by (2))

(by Proposition 2.3)

which completes the proof. Having shown that outer measure coincides with the natural concept of length for intervals, we now need to investigate its properties. The next theorem gives us an important technical tool which will be used in many proofs.

Theorem 2.5 Outer measure is countably subadditive, i.e. for any sequence of sets {En } m∗

∞ [

n=1

∞ X En ≤ m∗ (En ). n=1

(Note that both sides may be infinite here.)

Proof (a warm up) Let us prove first a simpler statement: m∗ (E1 ∪ E2 ) ≤ m∗ (E1 ) + m∗ (E2 ). Take an ε > 0 and we show an even easier inequality m∗ (E1 ∪ E2 ) ≤ m∗ (E1 ) + m∗ (E2 ) + ε. This is however sufficient because taking ε = n1 and letting n → ∞ we get what we need. So for any ε > 0 we find covering sequences (Ik1 )k≥1 of E1 and (Ik2 )k≥1 of E2 such that ∞ X ε l(Ik1 ) ≤ m∗ (E1 ) + , 2 k=1

∞ X k=1

hence, adding up, ∞ X

k=1

l(Ik1 ) +

l(Ik2 ) ≤ m∗ (E2 ) +

∞ X k=1

ε 2

l(Ik2 ) ≤ m∗ (E1 ) + m∗ (E2 ) + ε.

2. Measure

25

The sequence of intervals (I11 , I12 , I21 , I22 , I31 , I32 , . . .) covers E1 ∪ E2 hence m∗ (E1 ∪ E2 ) ≤

∞ X

l(Ik1 ) +

k=1

∞ X

l(Ik2 )

k=1

which combined with the previous inequality gives the result.

Proof (of the Theorem) If the right-hand side is infinite, then the inequality is of course true. So, suppose P∞ ∗ that n=1 m (En ) < ∞. For each given ε > 0 and n ≥ 1 find a covering n sequence (Ik )k≥1 of En with ∞ X k=1

l(Ikn ) ≤ m∗ (En ) +

ε . 2n

The iterated series converges: ∞ ∞ X X

n=1

k=1

∞ X l(Ikn ) ≤ m∗ (En ) + ε < ∞ n=1

and since all its terms are non-negative, ∞ ∞ X X

n=1

The system of intervals m∗

k=1

(Ikn )k,n≥1

∞ [

n=1

l(Ikn )

=

covers

∞ X

l(Ikn ).

n,k=1

S∞

n=1

En hence

∞ ∞ X X En ≤ l(Ikn ) ≤ m∗ (En ) + ε. n=1

n,k=1

To complete the proof we let ε → 0. A similar result is of course true for a finite family (En )m n=1 : m∗

m [

n=1

m X En ≤ m∗ (En ). n=1

It is a corollary to Theorem 2.5 with Ek = Ø for k > m.

Exercise 2.4 Prove that if m∗ (A) = 0 then for each B, m∗ (A ∪ B) = m∗ (B).

26


Hint Employ both monotonicity and subadditivity of outer measure.

Exercise 2.5 Prove that if m∗ (A∆B) = 0, then m∗ (A) = m∗ (B). Hint Note that A ⊆ B ∪ (A∆B). We conclude this section with a simple and intuitive property of outer measure. Note that the length of an interval does not change if we shift it along the real line: l([a, b]) = l([a + t, b + t]) = b − a for example. Since the outer measure is defined in terms of the lengths of intervals, it is natural to expect it to share this property. For A ⊂ R and t ∈ R we put A + t = {a + t : a ∈ A}.

Proposition 2.6 Outer measure is translation invariant, i.e. m∗ (A) = m∗ (A + t) for each A and t. Hint Combine two facts: the length of interval does not change when the interval is shifted and outer measure is determined by the length of the coverings.

2.3 Lebesgue measurable sets and Lebesgue measure With outer measure, subadditivity (as in Theorem 2.5 is as far as we can get. We wish, however, to ensure that if sets (En ) are pairwise disjoint (i.e. Ei ∩ Ej = Ø if i 6= j), then the inequality in Theorem 2.5 becomes an equality. It turns out that this will not in general be true for outer measure, although examples where it fails are quite difficult to construct (we give such examples in the Appendix). But our wish is an entirely reasonable one: any ‘length function’ should at least be finitely additive, since decomposing a set into finitely many disjoint pieces should not alter its length. Moreover, since we constructed our length function via approximation of complicated sets by ‘simpler’ sets (i.e. intervals) it seems fair to demand a continuity property: if pairwise disjoint (En ) have union E,

2. Measure

27

S then the lengths of the sets Bn = E \ nk=1 Ek may be expected to decrease to 0 as n → ∞. Combining this with finite additivity leads quite naturally to the demand that ‘length’ should be countably additive, i.e. that ! ∞ ∞ [ X ∗ m En = m∗ (En ) when Ei ∩ Ej = Ø for i 6= j. n=1

n=1

We therefore turn to the task of finding the class of sets in R which have this property. It turns out that it is also the key property of the abstract concept of measure, and we will use it to provide mathematical foundations for probability. In order to define the ‘good’ sets which have this property, it also seems plausible that such a set should apportion the outer measure of every set in R properly, as we state in Definition 2.3 below. Remarkably, this simple demand will suffice to guarantee that our ‘good’ sets have all the properties we demand of them!

Definition 2.3 A set E ⊆ R is (Lebesgue) measurable if for every set A ⊆ R we have m∗ (A) = m∗ (A ∩ E) + m∗ (A ∩ E c )

(2.6)

where E c = R\E, and we write E ∈ M. We obviously have A = (A ∩ E) ∪ (A ∩ E c ) hence by Theorem 2.5 we have m∗ (A) ≤ m∗ (A ∩ E) + m∗ (A ∩ E c ) for any A and E. So our future task of verifying (2.6) has simplified: E ∈ M if and only if the following inequality holds m∗ (A) ≥ m∗ (A ∩ E) + m∗ (A ∩ E c ) for all A ⊆ R.

(2.7)

Now we give some examples of measurable sets.

Theorem 2.7 (i) Any null set is measurable. (ii) Any interval is measurable.

Proof (i) If N is a null set, then (Proposition 2.2) m∗ (N ) = 0. So for any A ⊂ R we have m∗ (A ∩ N ) ≤ m∗ (N ) = 0 since A ∩ N ⊆ N

28


m∗ (A ∩ N c ) ≤ m∗ (A)

since A ∩ N c ⊆ A

and adding together we have proved (2.7). (ii) Let E = I be an interval. Suppose, for example, that I = [a, b]. Take any A ⊆ R and ε > 0. Find a covering of A with ∗

m (A) ≤

∞ X

n=1

l(In ) ≤ m∗ (A) + ε.

Clearly the intervals In0 = In ∩ [a, b] cover A ∩ [a, b] hence m∗ (A ∩ [a, b]) ≤

∞ X

l(In0 ).

n=1

The intervals In00 = In ∩ (−∞, a), In000 = In ∩ (b, +∞) cover A ∩ [a, b]c so ∗

c

m (A ∩ [a, b] ) ≤

∞ X

l(In00 )

n=1

+

∞ X

l(In000 ).

n=1

Putting the above three inequalities together we obtain (2.7). If I is unbounded, I = [a, ∞) say, then the proof is even simpler since it is sufficient to consider In0 = In ∩ [a, ∞) and In00 = In ∩ (−∞, a). The fundamental properties of the class M of all Lebesgue-measurable subsets of R can now be proved. They fall into two categories: first we show that certain set operations on sets in M again produce sets in M (these are what we call ‘closure properties’) and second we prove that for sets in M the outer measure m∗ has the property of countable additivity announced above.

Theorem 2.8 (i) R ∈ M,

(ii) if E ∈ M then E c ∈ M,

S∞ (iii) if En ∈ M for all n = 1, 2, . . . then n=1 En ∈ M. Moreover, if En ∈ M, n = 1, 2, . . . and Ej ∩ Ek = Ø for j 6= k, then m

∗

∞ [

n=1

En =

∞ X

n=1

m∗ (En ).

(2.8)

2. Measure

29

Remark 2.1 This result is the most important theorem in this chapter and provides the basis for all that follows. It also allows us to give names to the quantities under discussion. Conditions (i)–(iii) mean that M is a σ-field. In other words, we say that a family of sets is a σ-field if it contains the base set and is closed under complements and countable unions. A [0, ∞]-valued function defined on a σfield is called a measure if it satisfies (2.8) for pairwise disjoint sets, i.e. it is countably additive. An alternative, rather more abstract and general, approach to measure theory is to begin with the above properties as axioms, i.e. to call a triple (Ω, F, µ) a measure space if Ω is an abstractly given set, F is a σ-field of subsets of Ω, and µ : F 7→ [0, ∞] is a function satisfying (2.8) (with µ instead of m∗ ). The task of defining Lebesgue measure on R then becomes that of verifying, with M and m = m∗ on M defined as above, that the triple (R, M, m) satisfies these axioms, i.e. becomes a measure space. Although the requirements of probability theory will mean that we have to consider such general measure spaces in due course, we have chosen our more concrete approach to the fundamental example of Lebesgue measure in order to demonstrate how this important measure space arises quite naturally from considerations of the ‘lengths’ of sets in R, and leads to a theory of integration which greatly extends that of Riemann. It is also sufficient to allow us to develop most of the important examples of probability distributions.

Proof (of the Theorem) (i) Let A ⊆ R. Note that A ∩ R = A, Rc = Ø, so that A ∩ Rc = Ø. Now (2.6) reads m∗ (A) = m∗ (A) + m∗ (Ø) and is of course true since m∗ (Ø) = 0. (ii) Suppose E ∈ M and take any A ⊆ R. We have to show (2.6) for E c , i.e. m∗ (A) = m∗ (A ∩ E c ) + m∗ (A ∩ (E c )c ) but since (E c )c = E this reduces to the condition for E which holds by hypothesis. We split the proof of (iii) into several steps. But first: A warm up. Suppose that E1 ∩ E2 = Ø, E1 , E2 ∈ M. We shall show that E1 ∪ E2 ∈ M and m∗ (E1 ∪ E2 ) = m∗ (E1 ) + m∗ (E2 ).

30


Let A ⊆ R. We have the condition for E1 : m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E1c ).

(2.9)

Now, apply (2.6) for E2 with A ∩ E1c in place of A: m∗ (A ∩ E1c ) = m∗ ((A ∩ E1c ) ∩ E2 ) + m∗ ((A ∩ E1c ) ∩ E2c ). = m∗ (A ∩ (E1c ∩ E2 )) + m∗ (A ∩ (E1c ∩ E2c ))

(the situation is depicted in Figure 2.2).

Figure 2.2 The sets A, E1 , E2 Since E1 and E2 are disjoint, E1c ∩ E2 = E2 . By de Morgan’s law E1c ∩ E2c = (E1 ∪ E2 )c . We substitute and we have m∗ (A ∩ E1c ) = m∗ (A ∩ E2 ) + m∗ (A ∩ (E1 ∪ E2 )c ). Substituting this into (2.9) we get m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E2 ) + m∗ (A ∩ (E1 ∪ E2 )c ).

(2.10)

Now by the subadditivity property of m∗ we have m∗ (A ∩ E1 ) + m∗ (A ∩ E2 ) ≥ m∗ (A ∩ E1 ) ∪ (A ∩ E2 ) = m∗ (A ∩ (E1 ∪ E2 ))

so (2.10) gives m∗ (A) ≥ m∗ (A ∩ (E1 ∪ E2 )) + m∗ (A ∩ (E1 ∪ E2 )c ) which is sufficient for E1 ∪ E2 to belong to M (the inverse inequality is always true, as observed before (2.7)). Finally, put A = E1 ∪ E2 in (2.10) to get m∗ (E1 ∪ E2 ) = m∗ (E1 ) + m∗ (E2 ), which completes the argument. We return to the proof of the theorem.

2. Measure

31

Proof Step 1. If pairwise disjoint Ek , k = 1, 2, . . ., are in M then their union is in M and (2.8) holds. We begin as in the proof of Warm up and we have m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E1c ) m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E2 ) + m∗ (A ∩ (E1 ∪ E2 )c ) (see (2.10)) and after n steps we expect n X

m∗ (A) =

k=1

m∗ (A ∩ Ek ) + m∗ A ∩

n [

Ek

k=1

c

.

(2.11)

Let us demonstrate this by induction. The case n = 1 is the first line above. Suppose that n−1 X

m∗ (A) =

k=1

m∗ (A ∩ Ek ) + m∗ A ∩

Since En ∈ M, we may apply (2.6) with A ∩ m∗ (A ∩ (

n−1 [ k=1

Ek )c ) = m∗ (A ∩ (

n−1 [ k=1

Sn−1 k=1

n−1 [

Ek

k=1

Ek

c

Ek )c ∩ En ) + m∗ (A ∩ (

c

.

(2.12)

in place of A:

n−1 [ k=1

Ek )c ∩ Enc ). (2.13)

Now we make the same observations as in the Warm up: n−1 [

Ek

k=1

n−1 [

Ek

k=1

c

c

∩ En = E n

∩ Enc =

n [

k=1

Inserting these into (2.13) we get m∗ (A ∩

n−1 [ k=1

(Ei are pairwise disjoint),

Ek

c

(by de Morgan’s law).

c Ek ) = m∗ (A ∩ En ) + m∗ (A ∩

n [

k=1

c Ek ),

and inserting this into the induction hypothesis (2.12) we get m∗ (A) =

n−1 X k=1

m∗ (A ∩ Ek ) + m∗ (A ∩ En ) + m∗ A ∩

n [

k=1

Ek

c

as required to complete the induction step. Thus (2.11) holds for all n by induction.

32


As will be seen at the next step the fact that Ek are pairwise disjoint is not necessary in order to ensure that their union belongs to M. However, with this assumption we have equality in (2.11) which does not hold otherwise. This equality will allow us to prove countable additivity (2.8). Since

n [

Ek

k=1

c

⊇

∞ [

Ek

k=1

c

,

from (2.11) by monotonicity (Proposition 2.3) we get ∗

m (A) ≥

n X

k=1

∗

∗

m (A ∩ Ek ) + m A ∩

∞ [

Ek

k=1

c

.

The inequality remains true after we pass to the limit n → ∞: m∗ (A) ≥

∞ X

k=1

m∗ (A ∩ Ek ) + m∗ A ∩

∞ [

Ek

k=1

c

.

(2.14)

c

(2.15)

c

(2.16)

By countable subadditivity (Theorem 2.5) ∞ X k=1

m∗ (A ∩ Ek ) ≥ m∗ A ∩

and so m∗ (A) ≥ m∗ A ∩

∞ [

k=1

∞ [

k=1

Ek + m ∗ A ∩

Ek

∞ [

Ek

k=1

S as required. So we have shown that ∞ k=1 Ek ∈ M and hence the two sides of (2.15) are equal. The right hand side of (2.14) is squeezed between the left and right of (2.15) which yields m∗ (A) =

∞ X

k=1

m∗ (A ∩ Ek ) + m∗ A ∩

∞ [

Ek

k=1

.

The equality here is a consequence of the assumption that Ek are pairwise S∞ disjoint. It holds for any set A so we may insert A = j=1 Ej . The last term S∞ on the right is zero because we have m∗ (Ø). Next ( j=1 Ej ) ∩ En = En and so we have (2.8). Step 2. If E1 , E2 ∈ M, then E1 ∪ E2 ∈ M (not necessarily disjoint). Again we begin as in the Warm up: m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E1c ).

(2.17)

2. Measure

33

Next, applying (2.6) to E2 with A ∩ E1c in place of A we get m∗ (A ∩ E1c ) = m∗ (A ∩ E1c ∩ E2 ) + m∗ (A ∩ E1c ∩ E2c ). We insert this into (2.17) to get m∗ (A) = m∗ (A ∩ E1 ) + m∗ (A ∩ E1c ∩ E2 ) + m∗ (A ∩ E1c ∩ E2c ).

(2.18)

By de Morgan’s law, E1c ∩ E2c = (E1 ∪ E2 )c so (as before) m∗ (A ∩ E1c ∩ E2c ) = m∗ (A ∩ (E1 ∪ E2 )c ).

(2.19)

By subadditivity of m∗ we have m∗ (A ∩ E1 ) + m∗ (A ∩ E1c ∩ E2 ) ≥ m∗ (A ∩ (E1 ∪ E2 )).

(2.20)

Inserting (2.19) and (2.20) into (2.18) we get m∗ (A) ≥ m∗ (A ∩ (E1 ∪ E2 )) + m∗ (A ∩ (E1 ∪ E2 )c ) as required. Step 3. If Ek ∈ M, k = 1, . . . , n, then E1 ∪ . . . ∪ En ∈ M (not necessarily disjoint). We argue by induction. There is nothing to prove for n = 1. Suppose the claim is true for n − 1. Then E1 ∪ . . . ∪ En = (E1 ∪ . . . ∪ En−1 ) ∪ En so that the result follows from Step 2. Step 4. If E1 , E2 ∈ M, then E1 ∩ E2 ∈ M. We have E1c , E2c ∈ M by (ii), E1c ∪ E2c ∈ M by Step 2, (E1c ∪ E2c )c ∈ M by (ii) again, but by de Morgan’s law the last set is equal to E1 ∩ E2 . Step 5. The general case: if E1 , E2 , . . . are in M, then so is

S∞

k=1

Ek

Let Ek ∈ M, k = 1, 2, . . .. We define an auxiliary sequence of pairwise disjoint sets Fk with the same union as Ek : F1 = E 1 F2 = E2 \E1 = E2 ∩ E1c

F3 = E3 \(E1 ∪ E2 ) = E3 ∩ (E1 ∪ E2 )c ...

Fk = Ek \(E1 ∪ . . . ∪ Ek−1 ) = Ek ∩ (E1 ∪ . . . ∪ Ek−1 )c ,

34


E1

E2 \ E1

E3 \ (E1 È E2 ) E4 \ (E1 ÈE2 È E 3 )

Figure 2.3 The sets Fk see Figure 2.3. By Steps 3 and 4 we know that all Fk are in M. By the very construction they are pairwise disjoint so by Step 1 their union is in M. We shall show that ∞ [

Fk =

k=1

∞ [

Ek .

k=1

This will complete the proof since the latter is now in M. The inclusion ∞ [

k=1

Fk ⊆

∞ [

Ek

k=1

is obvious since for each k, Fk ⊆ Ek by definition. For the inverse let a ∈ S∞ k=1 Ek . Put S = {n ∈ N : a ∈ En } which is non-empty since a belongs to the union. Let n0 = min S ∈ S. If n0 = 1, then a ∈ E1 = F1 . Suppose n0 > 1. So a ∈ En0 and, by the definition of n0 , a ∈ / E1 , . . . , a ∈ / En0 −1 . By the definition S∞ of Fn0 this means that a ∈ Fn0 so a is in k=1 Fk . Using de Morgan’s laws you should easily verify an additional property of

M.

Proposition 2.9 If Ek ∈ M, k = 1, 2, . . ., then E=

∞ \

k=1

Ek ∈ M.

We can therefore summarize the properties of the family M of Lebesgue measurable sets as follows: 1. M is closed under countable unions, countable intersections, and complements. It contains intervals and all null sets.

2. Measure

35

Definition 2.4 We shall write m(E) instead of m∗ (E) for any E in M and call m(E) the Lebesgue measure of the set E. Thus Theorems 2.8 and 2.4 now read as follows, and describe the construction which we have laboured so hard to establish: 1. Lebesgue measure m : M → [0, ∞] is a countably additive set function defined on the σ-field M of measurable sets. Lebesgue measure of an interval is equal to its length. Lebesgue measure of a null set is zero.

2.4 Basic properties of Lebesgue measure Since Lebesgue measure is nothing else than the outer measure restricted to a special class of sets, some properties of the outer measure are automatically inherited by Lebesgue measure:

Proposition 2.10 Suppose that A, B ∈ M.

(i) If A ⊂ B then m(A) ≤ m(B).

(ii) If A ⊂ B and m(A) is finite then m(B \ A) = m(B) − m(A).

(iii) m is translation invariant.

Since Ø ∈ M we can take Ei = Ø for all i > n in (2.8) to conclude that Lebesgue measure is additive: if Ei ∈ M are pairwise disjoint, then m(

n [

i=1

Ei ) =

n X

m(Ei ).

i=1

Exercise 2.6 Find a formula describing m(A ∪ B) and m(A ∪ B ∪ C) in terms of measures of the individual sets and their intersections (we do not assume that the sets are pairwise disjoint). Recalling that the symmetric difference A∆B of two sets is defined by A∆B = (A \ B) ∪ (B \ A) the following result is also easy to check:

36


Proposition 2.11 If A ∈ M, m(A∆B) = 0, then B ∈ M and m(A) = m(B). Hint Recall that null sets belong to M and that subsets of null sets are null. As we noted in Chapter 1, every open set in R can be expressed as the union of a countable number of open intervals. This ensures that open sets in R are Lebesgue-measurable, since M contains intervals and is closed under countable unions. We can approximate the Lebesgue measure of any A ∈ M from above by the measures of a sequence of open sets containing A. This is clear from the following result:

Theorem 2.12 (i) For any ε > 0, A ⊂ R we can find an open set O such that A ⊂ O,

m(O) ≤ m∗ (A) + ε.

Consequently, for any E ∈ M we can find an open set O containing E such that m(O \ E) < ε. (ii) For any A ⊂ R we can find a sequence of open sets On such that \ \ A⊂ On , m( On ) = m∗ (A). n

n

Proof (i) By definition of m∗ (A) we can find a sequence (In ) of intervals with A ⊂ S P∞ ε ∗ n In and n=1 l(In ) − 2 ≤ m (A). Each In is contained in an open interval whose length is very close to that of In ; if the left and right endpoints of In are S ε ε , bn + 2n+2 ). Set O = n Jn , which an and bn respectively let Jn = (an − 2n+2 is open. Then A ⊂ O and m(O) ≤

∞ X

n=1

l(Jn ) ≤

∞ X

n=1

l(In ) +

ε ≤ m∗ (A) + ε. 2

When m(E) < ∞ the final statement follows at once from (ii) in Proposition 2.10, since then m(O \ E) = m(O) − m(E) ≤ ε. When m(E) = ∞ we S first write R as a countable union of finite intervals: R = n (−n, n). Now

2. Measure

37

En = E ∩ (−n, n) has finite measure, so we can find an open On ⊃ En with S m(On \ En ) ≤ 2εn . The set O = n On is open and contains E. Now [ [ [ O \ E = ( On ) \ ( En ) ⊂ (On \ En ) n

so that m(O \ E) ≤

P

n

n

n

m(On \ En ) ≤ ε as required.

T (ii) In (i) use ε = n1 and let On be the open set so obtained. With E = n On we obtain a measurable set containing A such that m(E) < m(On ) ≤ m∗ (A)+ n1 for each n, hence the result follows.

Remark 2.2 Theorem 2.12 shows how the freedom of movement allowed by the closure properties of the σ-field M can be exploited by producing, for any set A ⊂ R, a measurable set O ⊃ A which is obtained from open intervals with two operations (countable unions followed by countable intersections) and whose measure equals the outer measure of A. Finally we show that monotone sequences of measurable sets behave as one would expect with respect to m.

Theorem 2.13 Suppose that An ∈ M for all n ≥ 1. Then we have:

(i) if An ⊂ An+1 for all n, then [ m( An ) = lim m(An ), n

n→∞

(ii) if An ⊃ An+1 for all n and m(A1 ) < ∞, then \ m( An ) = lim m(An ). n

n→∞

38


Proof (i) Let B1 = A1 , Bi = Ai − Ai−1 for i > 1. Then Bi ∈ M are pairwise disjoint, so that [ [ m( Ai ) = m( Bi ) i

S∞

i=1

Bi =

S∞

i=1

Ai and the

i

=

∞ X

m(Bi ) (by countable additivity)

i=1

= lim

n→∞

n X

= lim m( n→∞

m(Bi )

i=1

n [

Bi ) (by additivity)

n=1

= lim m(An ), n→∞

since An =

Sn

i=1

Bi by construction – see Figure 2.4.

Figure 2.4 Sets An , Bn (ii) A1 \ A1 = Ø ⊂ A1 \ A2 ⊂ . . . ⊂ A1 \ An ⊂ . . . for all n, so that by (i) [ m( (A1 \ An )) = lim m(A1 \ An ) n

n→∞

and since m(A1 ) is finite, m(A1 \ An ) = m(A1 ) − m(An ). On the other hand, S T n (A1 \ An ) = A1 \ n An , so that [ \ m( (A1 \ An )) = m(A1 ) − m( An ) = m(A1 ) − lim m(An ). n

The result follows.

n

n→∞

2. Measure

39

Remark 2.3 The proof of Theorem 2.13 simply relies on the countable additivity of m and on the definition of the sum of a series in [0, ∞], i.e. that ∞ X

m(Ai ) = lim

n→∞

i=1

n X

m(Ai ).

i=1

Consequently the result is true, not only for the set function m we have constructed on M, but for any countably additive set function defined on a σ-field. It also leads us to the following claim, which, though we consider it here only for m, actually characterizes countably additive set functions.

Theorem 2.14 The set function m satisfies: (i) m is finitely additive, i.e. for pairwise disjoint sets (Ai ) we have m(

n [

Ai ) =

i=1

n X

m(Ai )

i=1

for each n; (ii) m is continuous at Ø, i.e. if (Bn ) decrease to Ø, then m(Bn ) decreases to 0.

Proof To prove this claim, recall that m : M 7→ [0, ∞] is countably additive. This implies (i), as we have already seen. To prove (ii), consider a sequence (Bn ) in M which decreases to Ø. Then An = Bn \ Bn+1 defines a disjoint sequence in S M, and n An = B1 . We may assume that B1 is bounded, so that m(Bn ) is finite for all n, so that, by Proposition 2.10 (ii), m(An ) = m(Bn )−m(Bn+1 ) ≥ 0 and hence we have m(B1 ) =

∞ X

m(An )

n=1

= lim

k→∞

k X

n=1

[m(Bn ) − m(Bn+1 )]

= m(B1 ) − lim m(Bn ) n→∞

which shows that m(Bn ) → 0, as required.

40


2.5 Borel sets The definition of M does not easily lend itself to verification that a particular set belongs to M; in our proofs we have had to work quite hard to show that M is closed under various operations. It is therefore useful to add another construction to our armoury; one which shows more directly how open sets (and indeed open intervals) and the structure of σ-fields lie at the heart of many of the concepts we have developed. We begin with an auxiliary construction enabling us to produce new σ-fields.

Theorem 2.15 The intersection of a family of σ-fields is a σ-field.

Proof Let Fα be σ-fields for α ∈ Λ (the index set Λ can be arbitrary). Put \ F= Fα . α∈Λ

We verify the conditions of the definition. 1. R ∈ Fα for all α ∈ Λ so R ∈ F. 2. If E ∈ F, then E ∈ Fα for all α ∈ Λ. Since the Fα are σ-fields, E c ∈ Fα and so E c ∈ F. S∞ 3. If Ek ∈ F for k = 1, 2, . . . , then Ek ∈ Fα , all α, k, hence k=1 Ek ∈ Fα , S∞ all α, and so k=1 Ek ∈ F.

Definition 2.5 Put B=

\

{F : F is a σ-field containing all intervals}.

We say that B is the σ-field generated by all intervals and we call the elements of B Borel sets (after Emile Borel 1871–1956). It is obviously the smallest σfield containing all intervals. In general, we say that G is the σ-field generated T by a family of sets A if G = {F : F is a σ-field such that F ⊃ A}.

Example 2.1 (Borel sets) The following examples illustrate how the closure properties of the σ-field B may be used to verify that most familiar sets in R belong to B.

2. Measure

41

(i) By construction, all intervals belong to B, and since B is a σ-field, all open sets must belong to B, as any open set is a countable union of (open) intervals. (ii) Countable sets are Borel sets, since each is a countable union of closed intervals of the form [a, a]; in particular N and Q are Borel sets. Hence, as the complement of a Borel set, the set of irrational numbers is also Borel. Similarly, finite and cofinite sets are Borel sets. The definition of B is also very flexible – as long as we start with all intervals of a particular type, these collections generate the same Borel σ-field:

Theorem 2.16 If instead of the family of all intervals we take all open intervals, all closed intervals, all intervals of the form (a, ∞) (or of the form [a, ∞), (−∞, b), or (−∞, b]), all open sets, or all closed sets, then the σ-field generated by them is the same as B.

Proof Consider for example the σ-field generated by the family of open intervals OI and denote it by C: \ C = {F ⊃ OI, F is a σ-field}.

We have to show that B = C. Since open intervals are intervals, OI ⊂ I (the family of all intervals), then {F ⊃ I} ⊂ {F ⊃ OI}

i.e. the collection of all σ-fields F which contain I is smaller than the collection of all σ-fields which contain the smaller family OI, since it is a more demanding requirement to contain a bigger family, so there are fewer such objects. The inclusion is reversed after we take the intersection on both sides, thus C ⊂ B (the intersection of a smaller family is bigger, as the requirement of belonging to each of its members is a less stringent one). We shall show that C contains all intervals. This will be sufficient, since B is the intersection of such σ-fields, so it is contained in each, so B ⊂ C. To this end consider intervals [a, b), [a, b], (a, b] (the intervals of the form (a, b) are in C by definition): [a, b) =

∞ \

n=1

(a −

1 , b), n

42


[a, b] =

∞ \

n=1

(a, b] =

(a −

∞ \

n=1

1 1 , b + ), n n

(a, b +

1 ). n

C as a σ-field is closed with respect to countable intersection, so it contains the sets on the right. The argument for unbounded intervals is similar. The proof is complete.

Exercise 2.7 Show that the family of intervals of the form (a, b] also generates the σ-field of Borel sets. Show that the same is true for the family of all intervals [a, b).

Remark 2.4 Since M is a σ-field containing all intervals, and B is the smallest such σ-field, we have the inclusion B ⊂ M, i.e. every Borel set in R is Lebesgue-measurable. The question therefore arises whether these σ-fields might be the same. In fact the inclusion is proper. It is not altogether straightforward to construct a set in M \ B, and we shall not attempt this here (but see the Appendix). However, by Theorem 2.12 (ii), given any E ∈ M we can find a Borel set B ⊃ E of the form B = ∩n On , where the (On ) are open sets, and such that m(E) = m(B). In particular, m(B∆E) = m(B \ E) = 0. Hence m cannot distinguish between the measurable set E and the Borel set B we have constructed. Thus, given a Lebesgue-measurable set E we can find a Borel set B such that their symmetric difference E∆B is a null set. Now we know that E∆B ∈ M, and it is obvious that subsets of null sets are also null, and hence in M. However, we cannot conclude that every null set will be a Borel set (if B did contain all null sets then by Theorem 2.12 (ii) we would have B = M), and this points to an ‘incompleteness’ in B which explains why, even if we begin by defining m on intervals and then extend the definition to Borel sets, we would also need to extend it further in order to be able to identify precisely which sets are ‘negligible’ for our purposes. On the other hand, extension of the measure m to the σ-field M will suffice, since M does contain all m-null sets and all subsets of null sets also belong to M.

2. Measure

43

We show that M is the smallest σ-field on R with this property, and we say that M is the completion of B relative to m and (R, M, m) is complete (whereas the measure space (R, B, m) is not complete). More precisely, a measure space (X, F, µ) is complete if for all F ∈ F with µ(F ) = 0, for all N ⊂ F we have N ∈ F (and so µ(N ) = 0). The completion of a σ-field G, relative to a given measure µ, is defined as the smallest σ-field F containing G such that, if N ⊂ G ∈ G and µ(G) = 0, then N ∈ F.

Proposition 2.17 The completion of G is of the form {G ∪ N : G ∈ F, N ⊂ F ∈ F with µ(F ) = 0}. This allows us to extend the measure µ uniquely to a measure µ ¯ on F by setting µ ¯(G ∪ N ) = µ(G) for G ∈ G.

Theorem 2.18 M is the completion of B.

Proof We show first that M contains all subsets of null sets in B: so let N ⊂ B ∈ B, B null, and suppose A ⊂ R. To show that N ∈ M we need to show that m∗ (A) ≥ m∗ (A ∩ N ) + m∗ (A ∩ N c ). First note that m∗ (A ∩ N ) ≤ m∗ (N ) ≤ m∗ (B) = 0. So it remains to show that m∗ (A) ≥ m∗ (A ∩ N c ) but this follows at once from monotonicity of m∗ . Thus we have shown that N ∈ M. Since M is a complete σ-field containing B, this means that M also contains the completion C of B. Finally, we show that M is the minimal such σ-field, i.e. that M ⊂ C: first T consider E ∈ M with m∗ (E) < ∞, and choose B = n On ∈ B as described above such that B ⊃ E, m(B) = m∗ (E). (We reserve the use of m for sets in B throughout this argument.) Consider N = B \ E, which is in M and has m∗ (N ) = 0, since m∗ is additive on M. By Theorem 2.12 (ii) we can find L ⊃ N , L ∈ B and m(L) = 0. In other words, N is a subset of a null set in B, and therefore E = B \ N belongs to the completion C of B. For E ∈ M with m∗ (E) = ∞, apply the above to En = E ∩ [−n, n] for each n ∈ N. Each m∗ (En ) is finite, so the En all

44


belong to C and hence so does their countable union E. Thus M ⊂ C and so they are equal. Despite these technical differences, measurable sets are never far from ‘nice’ sets, and, in addition to approximations from above by open sets, as observed in Theorem 2.12, we can approximate the measure of any E ∈ M from below by those of closed subsets.

Theorem 2.19 If E ∈ M then for given ε > 0 there exists a closed set F ⊂ E such that S m(E \ F ) < ε. Hence there exists B ⊂ E in the form B = n Fn , where all the Fn are closed sets, and m(E \ B) = 0.

Proof The complement E c is measurable and by Theorem 2.12 we can find an open set O containing E c such that m(O \ E c ) ≤ ε. But O \ E c = O ∩ E = E \ O c , and F = Oc is closed and contained in E. Hence this F is what we need. The final part is similar to Theorem 2.12 (ii), and the proof is left to the reader.

Exercise 2.8 Show that each of the following two statements is equivalent to saying that E ∈ M: (i) given ε > 0 there is an open set O ⊃ E with m∗ (O \ E) < ε, (ii) given ε > 0 there is a closed set F ⊂ E with m∗ (E \ F ) < ε.

Remark 2.5 The two statements in the above Exercise are the key to a considerable generalization, linking the ideas of measure theory to those of topology: A non-negative countably additive set function µ defined on B is called a regular Borel measure if for every Borel set B we have: µ(B) = inf{µ(O) : O open, O ⊃ B},

µ(B) = sup{µ(F ) : F closed, F ⊂ B}. In Theorems 2.12 and 2.19 we have verified these relations for Lebesgue measure. We shall consider other concrete examples of regular Borel measures later.

2. Measure

45

2.6 Probability The ideas which led to Lebesgue measure may be adapted to construct measures generally on arbitrary sets: any set Ω carrying an outer measure (i.e. a mapping from P (Ω) to [0, ∞] monotone and countably sub-additive) can be equipped with a measure µ defined on an appropriate σ-field F of its subsets. The resulting triple (Ω, F, µ) is then called a measure space, as observed in Remark 2.1. Note that in the construction of Lebesgue measure we only used the properties, not the particular form of the outer measure. For the present, however, we shall be content with noting simply how to restrict Lebesgue measure to any Lebesgue measurable subset B of R with m(B) > 0: Given Lebesgue measure m on the Lebesgue σ-field M let MB = {A ∩ B : A ∈ M} and for A ∈ MB write

mB (A) = m(A).

Proposition 2.20 (B, MB , mB ) is a complete measure space. Hint

S

i (Ai

∩ B) = (

S

i

Ai ) ∩ B and (A1 ∩ B) \ (A2 ∩ B) = (A1 \ A2 ) ∩ B.

We can finally state precisely what we mean by ‘selecting a number from [0,1] at random’: restrict Lebesgue measure m to the interval B = [0, 1] and consider the σ-field of M[0,1] of measurable subsets of [0, 1]. Then m[0,1] is a measure on M[0,1] with ‘total mass’ 1. Since all subintervals of [0,1] with the same length have equal measure, the ‘mass’ of m[0,1] is spread uniformly over 1 [0,1], so that, for example, the ‘probability’ of choosing a number from [0, 10 ) 6 7 1 is the same as that of choosing a number from [ 10 , 10 ), namely 10 . Thus all numerals are equally likely to appear as first digits of the decimal expansion of the chosen number. On the other hand, with this measure, the probability that the chosen number will be rational is 0, as is the probability of drawing an element of the Cantor set C. We now have the basis for some probability theory, although a general development still requires the extension of the concept of measure from R to abstract sets. Nonetheless the building blocks are already evident in the detailed development of the example of Lebesgue measure. The main idea in providing a

46


mathematical foundation for probability theory is to use the concept of measure to provide the mathematical model of the intuitive notion of probability. The distinguishing feature of probability is the concept of independence, which we introduce below. We begin by defining the general framework.

2.6.1 Probability space Definition 2.6 A probability space is a triple (Ω, F, P ) where Ω is an arbitrary set, F is a σ-field of subsets of Ω, and P is a measure on F such that P (Ω) = 1, called probability measure or briefly probability.

Remark 2.6 The original definition, given by Kolmogorov in 1932, is a variant of the above (see Theorem 2.14): (Ω, F, P ) is a probability space if (Ω, F) are given as above, and P is a finitely additive set function with P (Ø) = 0 and P (Ω) = 1 such that P (Bn ) & 0 whenever (Bn ) in F decreases to Ø.

Example 2.2 We see at once that Lebesgue measure restricted to [0, 1] is a probability measure. More generally: suppose we are given an arbitrary Lebesgue measurable 1 set Ω ⊂ R, with m(Ω) > 0. Then P = c·mΩ , where c = m(Ω) , and m = mΩ denotes the restriction of Lebesgue measure to measurable subsets of Ω, provides a probability measure on Ω, since P is complete and P (Ω) = 1. 1 For example, if Ω = [a, b], we obtain c = b−a , and P becomes the ‘uniform distribution’ over [a, b]. However, we can also use less familiar sets for our base 1 space; for example, Ω = [a, b] ∩ (R \ Q), c = b−a gives the same distribution over the irrationals in [a, b].

2.6.2 Events: conditioning and independence The word ‘event’ is used to indicate that something is happening. In probability a typical event is to draw elements from a set and then the event is concerned with the outcome belonging to a particular subset. So, as described above, if

2. Measure

47

Ω = [0, 1] we may be interested in the fact that a number drawn at random from [0, 1] belongs to some A ⊂ [0, 1]. We want to estimate the probability of this happening, and in the mathematical setup this is the number P (A), here m[0,1] (A). So it is natural to require that A should belong to M[0,1] , since these are the sets we may measure. By a slight abuse of the language, probabilists tend to identify the actual ‘event’ with the set A which features in the event. The next definition simply confirms this abuse of language.

Definition 2.7 Given a probability space (Ω, F, P ) we say that the elements of F are events. Suppose next that a number has been drawn from [0, 1] but has not been revealed yet. We would like to bet on it being in [0, 14 ] and we get a tip that it certainly belongs to [0, 12 ]. Clearly, given this ‘inside information’, the probability of success is now 21 rather than 14 . This motivates the following general definition.

Definition 2.8 Suppose that P (B) > 0. Then the number P (A|B) =

P (A ∩ B) P (B)

is called the conditional probability of A given B.

Proposition 2.21 The mapping A 7→ P (A|B) is countably additive on the σ-field FB . Hint Use the fact that A 7→ P (A ∩ B) is countably additive on F. A classical application of the conditional probability is the total probability formula which enables the computation of the probability of an event by means of conditional probabilities given some disjoint hypotheses:

Exercise 2.9 Prove that if Hi are pairwise disjoint events such that P (Hi ) 6= 0, then ∞ X P (A) = P (A|Hi )P (Hi ). i=1

S∞

i=1

Hi = Ω,

48


It is natural to say that the event A is independent of B if the fact that B takes place has no influence on the chances of A, i.e. P (A|B) = P (A). By definition of P (A|B) this immediately implies the relation P (A ∩ B) = P (A) · P (B) which is usually taken as the definition of independence. The advantage of this practice is that we may dispose of the assumption P (B) > 0.

Definition 2.9 The events A, B are independent if P (A ∩ B) = P (A) · P (B).

Exercise 2.10 Suppose that A and B are independent events. Show that Ac and B are also independent. The Exercise indicates that if A and B are independent events, then all elements of the σ-fields they generate are mutually independent, since these σfields are simply the collections FA = {Ø, A, Ac , Ω} and FB = {Ø, B, B c , Ω} respectively. This leads us to a natural extension of the definition: two σ-fields F1 and F2 are independent if for any choice of sets A1 ∈ F1 and A2 ∈ F2 we have P (A1 ∩ A2 ) = P (A1 )P (A2 ). However, the extension of these definitions to three or more events (or several σ-fields) needs a little care, as the following simple examples show:

Example 2.3 Let Ω = [0, 1], A = [0, 14 ] as before; then A is independent of B = [ 81 , 58 ] and of C = [ 81 , 83 ] ∪ [ 34 , 1]. In addition, B and C are independent. However, P (A ∩ B ∩ C) 6= P (A) · P (B) · P (C). Thus, given three events, the pairwise independence of each of the three possible pairs does not suffice for the extension of ‘independence’ to all three events. 1 On the other hand, with A = [0, 41 ], B = C = [0, 16 ] ∪ [ 14 , 11 16 ], (or alterna1 9 tively with C = [0, 16 ] ∪ [ 16 , 1]) P (A ∩ B ∩ C) = P (A) · P (B) · P (C)

(2.21)

2. Measure

49

but none of the pairs make independent events. This confirms further that we need to demand rather more if we wish to extend the above definition – pairwise independence is not enough, nor is (2.21); therefore we need to require both conditions to be satisfied together. Extending this to n events leads to:

Definition 2.10 The events A1 , . . . , An are independent if for all k ≤ n for each choice of k events, the probability of their intersection is the product of the probabilities. Again there is a powerful counterpart for σ-fields (which can be extended to sequences, and even arbitrary families):

Definition 2.11 The σ-fields F1 , F2 , ..., Fn defined on a given probability space (Ω, F, P ) are independent if, for all choices of distinct indices i1 , i2 , ..., ik from {1, 2, ..., n} and all choices of sets Fin ∈ Fin we have P (Fi1 ∩ Fi2 ∩ ... ∩ Fik ) = P (Fi1 ) · P (Fi2 ) · · · · · P (Fik ). The issue of independence will be revisited in the subsequent chapters where we develop some more tools to calculate probabilities

2.6.3 Applications to mathematical finance As indicated in the Preface, we will explore briefly how the ideas developed in each chapter can be applied in the rapidly growing field of mathematical finance. This is not intended as an introduction to this subject, but hopefully it will demonstrate how a consistent mathematical formulation can help to clarify ideas central to many disciplines. Readers who are unfamiliar with mathematical finance should consult texts such as [4], [5], [7] for definitions and a discussion of the main ideas of the subject. Probabilistic modelling in finance centres on the analysis of models for the evolution of the value of traded assets, such as stocks or bonds, and seeks to identify trends in their future behaviour. Much of the modern theory is concerned with evaluating derivative securities such as options, whose value is determined by the (random) future values of some underlying security, such as a stock.

50


We illustrate the above probability ideas on a classical model of stock prices, namely the binomial tree. This model is based on finitely many time instants at which the prices may change, and the changes are of a very simple nature. Suppose that the number of steps is N , denote the price at the k-th step by S(k), 0 ≤ k ≤ N. At each step the stock price changes in the following way: the price at a given step is the price at the previous step multiplied by U with probability p or D with probability q = 1 − p, where 0 < D < U. Therefore the final price depends on the sequence ω = (ω1 , ω2 , . . . , ωN ) where ωi = 1 indicates the application of the factor U or ωi = 0, which indicates application of the factor D. Such a sequence is called a path and we take Ω to consist of all possible paths. In other words, S(k) = S(0) × η(1) × · · · × η(k), where η(k) =

U with probability p, D with probability q.

Exercise 2.11 Suppose N = 5, U = 1.2, D = 0.9, and S(0) = 500. Find the number of all paths. How many paths lead to the price S(5) = 524.88? What is the probability that S(5) > 900 if the probability going up in a single step is 0.5? In general, the total number of paths is clearly 2N and at step k there are k + 1 possible prices. We construct a probability space by equipping Ω with the sigma field 2Ω of all subsets of Ω, and the probability defined on single-element sets by P ({ω}) = PN pk q n−k , where k = i=1 ωi . As time progresses we gather information about stock prices, or, what amounts to the same, about paths. This means, that having observed some prices the range of possible future developments is restricted. Our information increases with time and this idea can be captured by the following family of σ-fields. Fix m < n and define a σ-field Fm = {A : ω, ω 0 ∈ A =⇒ ω1 = ω10 , ω2 = 0 ω20 , . . . , ωm = ωm }. So all paths from a particular set A in this sigma field have identical initial segments while the remaining coordinates are arbitrary. Note that F0 = {Ω, Ø}, F1 = {A1 , Ac1 , Ω, Ø}, where A1 = {ω : ω1 = 1}, i.e. S(1) = S(0)U, and c A1 = {ω : ω1 = 0} i.e. S(1) = S(0)D.

2. Measure

51

Exercise 2.12 Prove that Fm has 22

m

elements.

Exercise 2.13 Prove that the sequence Fm is increasing. This sequence is an example of a filtration (the identifying features are that the sigma fields should be contained in F and form an increasing chain), a concept which we shall revisit later on. The consecutive choices of stock prices are closely related to coin tossing. Intuition tells us that the latter are independent. This can be formally seen by introducing another σ-field describing the fact that at a particular step we have a particular outcome. Suppose ω is such that ωk = 1. Then we can identify the set of all paths with this property Ak = {ω : ωk = 1} and extend to a σ-field: Gk = {Ak , Ack , Ω, Ø}. In fact, Ack = {ω : ωk = 0}.

Exercise 2.14 Prove that Gm and Gk are independent if m 6= k.

2.7 Proofs of propositions Proof (of Proposition 2.3) S If the intervals In cover B, then they also cover A: A ⊂ B ⊂ n In , hence ZB ⊂ ZA . The infimum of a larger set cannot be greater than the infimum of a smaller set (trivial illustration: inf{0, 1, 2} < inf{1, 2}, inf{0, 1, 2} = inf{0, 2}) hence the result.

Proof (of Proposition 2.6) If a system In of intervals covers A then the intervals In + t cover A + t. Conversely, if Jn cover A + t then Jn − t cover A. Moreover, the total length of a family of intervals does not change when we shift each by a number. So we have a one–one correspondence between the interval coverings of A and A + t and this correspondence preserves the total length of the covering. This implies that the sets ZA and ZA+t are the same so their infima are equal.

52


Proof (of Proposition 2.9) By de Morgan’s law ∞ \

k=1

Ekc

Ek =

∞ [

k=1

c Ekc .

By Theorem 2.8 (ii) all are in M, hence by (iii) the same can be said about S c the union ∞ E . Finally, by (ii) again, the complement of this union is in k=1 k T∞ M, and so the intersection k=1 Ek is in M.

Proof (of Proposition 2.10) (i) Proposition 2.3 tells us that the outer measure is monotone, but since m is just the restriction of m∗ to M, then the same is true for m: A ⊂ B implies m(A) = m∗ (A) ≤ m∗ (B) = m(B). (ii) We write B as a disjoint union B = A ∪ (B \ A) and then by additivity of m we have m(B) = m(A) + m(B \ A). Subtracting m(A) (here it is important that m(A) is finite) we get the result. (iii) Translation invariance of m follows at once from translation invariance of the outer measure in the same way as in (i) above.

Proof (of Proposition 2.11) The set A∆B is null hence so are its subsets A \ B and B \ A. Thus these sets are measurable, and so is A ∩ B = A \ (A \ B), and therefore also B = (A ∩ B) ∪ (B \ A) ∈ M. Now m(B) = m(A ∩ B) + m(B \ A) as the sets on the right are disjoint. But m(B \ A) = 0 = m(A \ B), so m(B) = m(A ∩ B) = m(A ∩ B) + m(A \ B) = m((A ∩ B) ∪ (A \ B)) = m(A).

Proof (of Proposition 2.17) The family G = {G∪N : G ∈ F, N ⊂ F ∈ F with µ(F ) = 0} contains the set X S S S since X ∈ F. If Gi ∪Ni ∈ G, Ni ⊂ Fi , µ(Fi ) = 0, then Gi ∪Ni = Gi ∪ Ni is in G since the first set on the right is in F and the second is a subset of a null set S Fi ∈ F. If G∪N ∈ G, N ⊂ F, then (G∪N )c = (G∪F )c ∪((F \N )∩Gc ).which is also in G Thus G is a σ-field. Consider any other σ-field H containing F and all subsets of null sets. Since H is closed with respect to the unions, it contains G and so G is the smallest σ-field with this property.

2. Measure

53

Proof (of Proposition 2.20) It follows at once from the definitions and the Hint that MB is a σ-field. To see that mB is a measure we check countable additivity: with Ci = Ai ∩ B pairwise disjoint in MB , we have X [ [ X X m(Ai ∩ B) = m(Ci ) = mB (Ci ). mB ( Ci ) = m( (Ai ∩ B)) = i

i

i

i

i

Therefore (B, MB , mB ) is a measure space. It is complete, since subsets of null sets contained in B are by definition mB -measurable.

Proof (of Proposition 2.21) Assume that An are measurable and pairwise disjoint. By the definition of conditional probability P(

∞ [

n=1

An |B) = = = =

∞ [ 1 P (( An ) ∩ B) P (B) n=1 ∞ [ 1 P( (An ∩ B)) P (B) n=1 ∞ 1 X P (An ∩ B) P (B) n=1

∞ X

n=1

P (An |B)

since An ∩ B are also pairwise disjoint and P is countably additive.

3 Measurable functions

3.1 The extended real line The length of R is unbounded above, i.e. ‘infinite’. To deal with this we defined Lebesgue measure for sets of infinite as well as finite measure. In order to handle functions between such sets comprehensively, it is convenient to allow functions which take infinite values: we take their range to be (part of) the ‘extended real line’ R = [−∞, ∞], obtained by adding the ‘points at infinity’ −∞ and +∞ to R. Arithmetic in this set needs a little care as already observed in Section 2.2: we assume that a+∞ = ∞ for all real a, a×∞ = ∞ for a > 0, a×∞ = −∞ for a < 0, ∞ × ∞ = ∞ and 0 × ∞ = 0, with similar definitions for −∞. These are all ‘obvious’ intuitively (except possibly 0 × ∞), and (as for measures) we avoid ever forming ‘sums’ of the form ∞+(−∞). With these assumptions ‘arithmetic works as before’.

3.2 Lebesgue-measurable functions The domain of the functions we shall be considering is usually R. Now we have the freedom of defining f only ‘up to null sets’: once we have shown two functions f and g to be equal on R \ E where E is some null set, then f = g for all practical purposes. To formalize this, we say that f has a property (P ) almost everywhere (a.e.) if f has this property at all points of its domain, except

55

56


possibly on some null set. For example, the function f (x) =

1 for x 6= 0 0 for x = 0

is almost everywhere continuous, since it is continuous on R \ {0}, and the exceptional set {0} is null. (Note: Probabilists tend to say ‘almost surely’ (a.s.) instead of ‘almost everywhere’ (a.e.) and we shall follow their lead in the sections devoted to probability.) The next definition will introduce the class of Lebesgue-measurable functions. The condition imposed on f : R → R will beR necessary (though not sufficient) to give meaning to the (Lebesgue) integral f dm. Let us first give some motivation. Integration is always concerned with the process of approximation. In the Riemann integral we split the interval I = [a, b], over which we integrate into small pieces In – again intervals. The simplest method of doing this is to divide the interval into N equal parts. Then we construct approximating sums by multiplying the lengths of the small intervals by certain numbers cn (related to the values of the function in question; for example cn = inf In f , cn = supIn f , or cn = f (x) for some x ∈ In ): N X

cn l(In ).

n=1

For large n this sum is close to the Riemann integral regularity of f ).

Rb a

f (x) dx (given some

Figure 3.1 Riemann vs. Lebesgue The approach to the Lebesgue integral is similar but there is a crucial difference. Instead of splitting the integration domain into small parts, we decompose

3. Measurable functions

57

the range of the function. Again, a simple way is to introduce short intervals Jn of equal length. To build the approximating sums we first take the inverse images of Jn by f , i.e. f −1 (Jn ). These may be complicated sets, not necessarily intervals. Here the theory of measure developed previously comes into its own. We are able to measure sets provided they are measurable, i.e. they are in M. Given that, we compute N X cn m(f −1 (Jn )) n=1

where cn ∈ Jn or cn = inf Jn , for example. The following definition guarantees that this procedure makes sense (though some extra care may be needed to arrive at a finite number as N → ∞).

Definition 3.1 Suppose that E is a measurable set. We say that a function f : E −→ R is (Lebesgue-)measurable if for any interval I ⊆ R f −1 (I) = {x ∈ R : f (x) ∈ I} ∈ M. In what follows, the term measurable (without qualification) will refer to Lebesgue-measurable functions. If all the sets f −1 (I) ∈ B, i.e. if they are Borel sets, we call f Borelmeasurable, or simply a Borel function. The underlying philosophy is one which is common for various mathematical notions: the inverse image of a nice set is nice. Remember continuous functions, for example, where the inverse image of any open set is required to be open. The actual meaning of the word nice depends on the particular branch of mathematics. In the above definitions, note that since B ⊂ M, every Borel function is (Lebesgue-)measurable.

Remark 3.1 The terminology is somewhat unfortunate. ‘Measurable’ objects should be measured (as with measurable sets). However, measurable functions will be integrated. This confusion stems from the fact that the word integrable which would probably fit best here, carries a more restricted meaning, as we shall see later. This terminology is widely accepted and we are not going to try to fight the whole world here. We give some equivalent formulations:

58


Theorem 3.1 The following conditions are equivalent (a) f is measurable, (b) for all a, f −1 ((a, ∞)) is measurable, (c) for all a, f −1 ([a, ∞)) is measurable,

(d) for all a, f −1 ((−∞, a)) is measurable, (e) for all a, f −1 ((−∞, a]) is measurable.

Proof Of course (a) implies any of the other conditions. We show that (b) implies (a). The proofs of the other implications are similar, and are left as exercises (which you should attempt). We have to show that for any interval I, f −1 (I) ∈ M. By (b) we have that for the particular case I = (a, ∞). Suppose I = (−∞, a]. Then f −1 ((−∞, a]) = f −1 (R \ (a, ∞)) = E \ f −1 ((a, ∞)) ∈ M

(3.1)

since both E and f −1 ((a, ∞)) are in M (we use the closure properties of M established before). Next f −1 ((−∞, b)) = f −1

∞ [

n=1

=

∞ [

(−∞, b −

1 ] n

1 f −1 (−∞, b − ] . n n=1

By (3.1), f −1 (−∞, b − n1 ] ∈ M and the same is true for the countable union. From this we can easily deduce that f −1 ([b, ∞)) ∈ M. Now let I = (a, b), and f −1 ((a, b)) = f −1 ((−∞, b) ∩ (a, ∞))

= f −1 ((−∞, b)) ∩ f −1 ((a, ∞))

is in M as the intersection of two elements of M. By the same reasoning M contains f −1 ([a, b]) = f −1 ((−∞, b] ∩ [a, ∞))

= f −1 ((−∞, b]) ∩ f −1 ([a, ∞))

and half-open intervals are handled similarly.


59

3.3 Examples The following simple results show that most of the functions encountered ‘in practice’ are measurable. (i) Constant functions are measurable. Let f (x) ≡ c. Then R if a < c f −1 ((a, ∞)) = Ø otherwise and in both cases we have measurable sets. (ii) Continuous functions are measurable. For we note that (a, ∞) is an open set and so is f −1 ((a, ∞)). As we know, all open sets are measurable.

(iii) Define the indicator function of a set A by 1 if x ∈ A 1A (x) = 0 otherwise. Then A∈M since

⇔

1A is measurable

  R if a < 0 1−1 ((a, ∞)) = A if 0 ≤ a < 1 A  Ø if a ≥ 1.

Exercise 3.1 Prove that every monotone function is measurable.

Exercise 3.2 Prove that if f is a measurable function, then the level set {x : f (x) = a} is measurable for every a ∈ R. Hint Don’t forget about the case when a is infinite!

Remark 3.2 In the Appendix, assuming the validity of the Axiom of Choice, we show that there are subsets of R which fail to be Lebesgue-measurable, and that there

60


are Lebesgue-measurable sets which are not Borel sets. Thus, if P(R) denotes the σ-field of all subsets of R, the following inclusions are strict B ⊂ M ⊂ P(R). These (rather esoteric) facts can be used, by considering the indicator functions of these sets, to construct examples of non-measurable functions and of measurable functions which are not Borel functions. While it is important to be aware of these distinctions in order to understand why these different concepts are introduced at all, such examples will not feature in the applications of the theory which we have in mind.

3.4 Properties The class of measurable functions is very rich, as the following results show.

Theorem 3.2 The set of real-valued measurable functions defined on E ∈ M is a vector space and closed under multiplication, i.e. if f and g are measurable functions then f + g, and f g are also measurable (in particular, if g is a constant function g ≡ c, cf is measurable for all real c).

Proof Fix measurable functions f, g : E → R. First consider f + g. Our goal is to show that for each a ∈ R, B = (f + g)−1 (−∞, a) = {t : f (t) + g(t) < a} ∈ M. Suppose that all the rationals are arranged in a sequence {qn }. Now B=

∞ [

n=1

{t : f (t) < qn , g(t) < a − qn }

– we decompose the half-plane below the line x + y = a into a countable union of unbounded ‘boxes’: {(x, y) : x < qn , y < a − qn }. Clearly {t : f (t) < qn , g(t) < a − qn } = {t : f (t) < qn } ∩ {t : g(t) < a − qn } is measurable as an intersection of measurable sets. Hence B ∈ M as a countable union of elements of M.


61

Figure 3.2 Boxes To deal with f g we adopt a slightly indirect approach in order to remain ‘one-dimensional’: first note that if g is measurable, then so is −g. Hence f −g = f +(−g) is measurable. Since f g = 41 {(f +g)2 −(f −g)2 }, it will suffice to prove that the square of a measurable function is measurable. So take a measurable h : E → R and consider {x ∈ E : h2 (x) > a}. For a < 0 this set is E ∈ M, and for a ≥ 0 √ √ {x : h2 (x) > a} = {x : h(x) > a} ∪ {x : h(x) < − a}. Both sets on the right are measurable, hence we have shown that h2 is measurable. Apply this with h = f + g and h = f − g respectively, to conclude that f g is measurable. It follows that cf is measurable for constant c, hence that the class of real-valued measurable functions forms a vector space under addition.

Remark 3.3 An elegant proof of the theorem is based on the following lemma, which will also be useful later. Its proof makes use of the simple topological fact that every open set in R2 decomposes into a countable union of rectangles, in precise analogy with open sets in R and intervals.

Lemma 3.3 Suppose that F : R×R → R is a continuous function. If f and g are measurable, then h(x) = F (f (x), g(x)) is also measurable.

62


It now suffices to take F (u, v) = u + v, F (u, v) = uv to obtain a second proof of Theorem 3.2.

Proof (of the Lemma) For any real a {x : h(x) > a} = {x : (f (x), g(x)) ∈ Ga } where Ga = {(u, v) : F (u, v) > a} = F −1 ((a, ∞)). Suppose for the moment that we have been lucky and Ga is a rectangle: Ga = (a1 , b1 ) × (c1 , d1 ).

Figure 3.3 The sets Ga It is clear from Figure 3.3 that {x : h(x) > a} = {x : f (x) ∈ (a1 , b1 ) and g(x) ∈ (c1 , d1 )}

= {x : f (x) ∈ (a1 , b1 )} ∩ {x : g(x) ∈ (c1 , d1 )}.

In general, we have to decompose the set Ga into a union of rectangles. The set Ga is an open subset of R × R since F is continuous. Hence it can be written as ∞ [ Ga = Rn n=1


63

where Rn are open rectangles Rn = (an , bn ) × (cn , dn ). So {x : h(x) > a} =

∞ [

n=1

{x : f (x) ∈ (an , bn )} ∩ {x : g(x) ∈ (cn , dn )}

is measurable due to the stability properties of M. A simple application of Theorem 3.2 is to consider the product f · 1A . If f is a measurable function, A is a measurable set, then f · 1A is measurable. This function is simply f on A and 0 outside A. Applying this to the set A = {x ∈ E : f (x) > 0} we see that the positive part f + of a measurable function is measurable: we have f (x) if f (x) > 0 + f (x) = 0 if f (x) ≤ 0. Similarly the negative part f − of f is measurable, since 0 if f (x) > 0 f − (x) = −f (x) if f (x) ≤ 0.

Proposition 3.4 Let E be a measurable subset of R. (i) f : E → R is measurable if and only if both f + and f − are measurable.

(ii) If f is measurable, then so is |f |; but the converse is false.

Hint Part (ii) requires the existence of non-measurable sets (as proved in the Appendix) not their particular form.

Exercise 3.3 Show that if f is measurable, then the truncation of f : a if f (x) > a f a (x) = f (x) if f (x) ≤ a is also measurable.

Exercise 3.4 Find a non-measurable f such that f 2 is measurable.

64


Passage to the limit does not destroy measurability – all the work needed was done when we established the stability properties of M!

Theorem 3.5 If {fn } is a sequence of measurable functions defined on the set E in R, then the following are measurable functions also: max fn , n≤k

min fn ,

sup fn ,

n≤k

n∈N

inf fn ,

lim sup fn ,

n∈N

n→∞

lim inf fn . n→∞

Proof It is sufficient to note that the following are measurable sets: {x : (max fn )(x) > a} = n≤k

{x : (min fn )(x) > a} = n≤k

{x : (sup fn )(x) > a} = n≥k

{x : ( inf fn )(x) ≥ a} = n≥k

k [

n=1 k \

n=1 ∞ [

{x : fn (x) > a}, {x : fn (x) > a},

n=k ∞ \

n=k

{x : fn (x) > a}, {x : fn (x) ≥ a}.

For the upper limit, by definition lim sup fn = inf { sup fm } n→∞

n≥1 m≥n

and the above relations show that hn = supm≥n fm is measurable, hence inf n≥1 hn (x) is measurable. The lower limit is done similarly.

Corollary 3.6 If a sequence fn of measurable functions converges (pointwise) then the limit is a measurable function.

Proof This is immediate since limn→∞ fn = lim supn→∞ fn which is measurable.


65

Remark 3.4 Note that Theorems 3.2 and 3.5 have counterparts for Borel functions, i.e. they remain valid upon replacing ‘measurable’ by ‘Borel’ throughout. Things are slightly more complicated when we consider the role of null sets. On the one hand, changing a function on a null set cannot destroy its measurability, i.e. any measurable function which is altered on a null set remains measurable. However, as not all null sets are Borel sets, we cannot conclude similarly for Borel sets, and thus the following results have no natural ‘Borel’ counterparts.

Theorem 3.7 If f : E → R is measurable, E ∈ M, g : E → R is arbitrary, and the set {x : f (x) = g(x)} is null, then g is measurable.

Proof Consider the difference d(x) = g(x) − f (x). It is zero except on a null set so a null set if a ≥ 0 {x : d(x) > a} = a full set if a < 0 where a full set is the complement of a null set. Both null and full sets are measurable hence d is a measurable function. Thus g = f +d is measurable.

Corollary 3.8 If (fn ) is a sequence of measurable functions and fn (x) → f (x) almost everywhere for x in E, then f is measurable.

Proof Let A be the null set such that fn (x) converges for all x ∈ E \ A. Then 1Ac fn converge everywhere to g = 1Ac f which is therefore measurable. But f = g almost everywhere, so f is also measurable.

Exercise 3.5 Let fn be a sequence of measurable functions. Show that the set E = {x : fn (x) converges} is measurable.

66


Since we are able to adjust a function f at will on a null set without altering its measurability properties, the following definition is a useful means of concentrating on the values of f that ‘really matter’ for integration theory, by identifying its bounds ‘outside null sets’:

Definition 3.2 Suppose f : E → R is measurable. The essential supremum ess sup f is defined as inf{z : f ≤ z a.e.} and the essential infimum ess inf f is sup{z : f ≥ z a.e.}. Note that ess sup f can be +∞. If ess sup f = −∞, then f = −∞ a.e. since by definition of ess sup, f ≤ −n a.e. for all n ≥ 1. Now if ess sup f is finite, and A = {x : ess sup f < f (x)}, define An for n ≥ 1 by 1 An = {x : ess sup f < f (x) − }. n S These are null sets, hence so is A = n An , and thus we have verified: f ≤ ess sup f a.e.

The following is now straightforward to prove.

Proposition 3.9 If f, g are measurable functions, then ess sup (f + g) ≤ ess sup f + ess sup g.

Exercise 3.6 Show that for measurable f , ess sup f ≤ sup f . Show that these quantities coincide when f is continuous.

3.5 Probability 3.5.1 Random variables In the special case of probability spaces we use the phrase random variable to mean a measurable function. That is, if (Ω, F, P ) is a probability space, then X : Ω → R is a random variable if for all a ∈ R the set X −1 ([a, ∞)) is in F: {ω ∈ Ω : X(ω) ≥ a} ∈ F.


67

In the case where Ω ⊂ R is a measurable set and F = B is the σ-field of Borel subsets of Ω, random variables are just Borel functions R → R. In applied probability, the set Ω represents the outcomes of a random experiment that can be observed by means of various measurements. These measurements assign numbers to outcomes and thus we arrive at the notion of random variable in a natural way. The condition imposed guarantees that questions of the following sort make sense: what is the probability that the value of the random variable lies within given limits?

3.5.2 Sigma fields generated by random variables As indicated before, the random variables we encounter will in fact be Borel measurable functions. The values of the random variable X will not lead us to non-Borel sets; in fact, they are likely to lead us to discuss much coarser distinctions between sets than are already available within the complexity of the Borel σ-field B. We should therefore be ready to consider different σ-fields contained within F. To be precise: The family of sets X −1 (B) = {S ⊂ F : S = X −1 (B) for some B ∈ B} is a σ-field. If X is a random variable, X −1 (B) ⊂ F but it may be a much smaller subset depending on the degree of sophistication of X. We denote this σ-field by FX and call it the σ-field generated by X. The simplest possible case is where X is constant, X ≡ a. The X −1 (B) is either Ω or Ø depending on whether a ∈ B or not and the σ-field generated is trivial: F = {Ø, Ω}. If X takes two values a 6= b, then FX contains four elements: FX = {Ø, Ω, X −1 ({a}), X −1 ({b})}. If X takes finitely many values, FX is finite. If X takes denumerably many values, FX is uncountable (it may be identified with the σ-field of all subsets of a countable set). We can see that the size of FX grows together with the level of complication of X.

Exercise 3.7 Show that FX is the smallest σ-field containing the inverse images X −1 (B) of all Borel sets B.

Exercise 3.8 Is the family of sets {X(A) : A ∈ F} a σ-field?

68


The notion of FX has the following interpretation. The values of the measurement X are all we can observe. From these we deduce some information on the level of complexity of the random experiment, that is the size of Ω and FX , and we can estimate the probabilities of the sets in FX by statistical methods. The σ-field generated represents the amount of information produced by the random variable. For example, suppose that a die is thrown and only 0 and 1 are reported depending on the number shown being odd or even. We will never distinguish this experiment from coin tossing. The information provided by the measurement is insufficient to explore the complexity of the experiment (which has six possible outcomes, here grouped together into two sets).

3.5.3 Probability distributions For any random variable X we can introduce a measure on the σ-field of Borel sets B by setting PX (B) = P (X −1 (B)). We call PX the probability distribution of the random variable X.

Theorem 3.10 The set function PX is countably additive.

Proof Given pairwise disjoint Borel sets Bi their inverse images X −1 (Bi ) are pairwise S S disjoint and X −1 ( i Bi ) = i X −1 (Bi ), so [ [ [ X PX ( Bi ) = P (X −1 ( Bi )) = P ( X −1 (Bi )) = P (X −1 (Bi )) i

i

=

X

i

i

PX (Bi )

i

as required. Thus (R, B, PX ) is a probability space. For this it is sufficient to note that PX (R) = P (Ω) = 1. We consider some simple examples. Suppose that X is constant, i.e. X ≡ a. Then we call PX the Dirac measure concentrated at a and denote by δa . Clearly 1 if a ∈ B δa (B) = 0 if a ∈ / B.


69

In particular, δa ({a}) = 1. If X takes 2 values: X(ω) =

a with probability p b with probability 1 − p,

then

and so

 1 if a, b ∈ B    p if a ∈ B, b ∈ /B PX (B) =  1 − p if b ∈ B, a ∈ /B   0 otherwise, PX (B) = pδa (B) + (1 − p)δb (B).

The distribution of a general discrete random variable (i.e. one which takes only finitely many different values, except possibly on some null set) is of the form: if the values of X are ai taken with probabilities pi > 0, i = 1, 2, . . . P pi = 1, then ∞ X PX (B) = pi δai (B). i=1

Classical examples are:

(i) the geometric distribution, where pi = (1 − q)q i for some q ∈ (0, 1),

(ii) the Poisson distribution where pi =

λi −λ . i! e

We shall not discuss the discrete case further since this is not our the primary goal in this text, and it is covered in many elementary texts on probability theory (such as [9]). Now consider the classical probability space with Ω = [0, 1], F = B, P = m|[0,1] – Lebesgue measure restricted to [0,1]. We can give examples of random variables given by explicit formulae. For instance, let X(ω) = aω + b. Then the image of [0, 1] is the interval [b, a + b] and PX = a1 m|[b,a+b] , i.e. for Borel B PX (B) =

m(B ∩ [b, a + b]) . a

Example 3.1 Suppose a car leaves city A at random between 12 am and 1 pm. It travels at 50 mph towards B which is 25 miles from A. What is the probability distribution of the distance between the car and B at 1 pm?

70


Clearly, this distance is 0 with probability 12 , i.e. if the car departs before 12.30. As a function of the starting time (represented as ω ∈ [0, 1]) the distance has the form 0 if ω ∈ [0, 21 ] X(ω) = 50ω − 25 if ω ∈ ( 12 , 1] 1 and PX = 12 P1 + 12 P2 where P1 = δ0 , P2 = 25 m[0,25] . In this example, therefore, PX is a combination of Dirac and Lebesgue measures.

In later chapters we shall explore more complicated forms of X and the corresponding distributions after developing further machinery needed to handle the computations.

3.5.4 Independence of random variables Definition 3.3 X, Y are independent if the σ-fields generated by them are independent. In other words, for any Borel sets B, C in R, P (X −1 (B) ∩ Y −1 (C)) = P (X −1 (B))P (Y −1 (C)).

Example 3.2 Let (Ω = [0, 1], M) be equipped with Lebesgue measure. Consider X = 1[0, 12 ] , Y = 1[ 41 , 34 ] . Then FX = {Ø, [0, 1], [0, 12 ], ( 12 , 1]}, FY = {Ø, [0, 1], [ 41 , 34 ], [0, 14 ) ∪ ( 34 , 1]} are clearly independent.

Example 3.3 Let Ω be as above and let X(ω) = ω, Y (ω) = 1 − ω. Then FX = FY = M. A σ-field cannot be independent with itself (unless it is trivial): Take A ∈ F and then independence requires P (A ∩ A) = P (A) × P (A) (the set A belongs to ‘both’ σ-fields), i.e. P (A) = P (A)2 which can happen only if either P (A) = 0 or P (A) = 1. So a σ-field independent with itself consists of sets of measure zero or one.


71

3.5.5 Applications to mathematical finance Consider a model of stock prices, discrete in time, i.e. assume that the stock prices are given by a sequence S(n) of random variables, n = 1, 2, . . . , N . If the length of one step is h, then we have the time horizon T = N h and we shall often write S(T ) instead of S(N ). An example of such a model is the binomial tree considered in the previous chapter. Recall that a European call option is the random variable of the form (S(N ) − K)+ (N is the exercise time, K is the strike price, S is the underlying asset). A natural generalisation of this is a random variable of the form f (S(N )) for some measurable function f : R → R. This random variable is of course measurable with respect to the σ-field generated by S(N ). This allows us to formulate a general definition:

Definition 3.4 A European derivative security (contingent claim) with the underlying asset represented by a sequence S(n) and exercise time N is a random variable X measurable with respect to the σ-field F generated by S(N ).

Proposition 3.11 A European derivative security X must be of the form X = f (S(N )) for some measurable real function f. The above definition is not sufficient for applications. For example, it does not cover one of the basic derivative instruments, namely futures. Recall that a holder of the futures contract has the right to receive (or an obligation to pay in case of negative values) a certain sequence (X(1), . . . , X(N )) of cash payments depending on the values of the underlying security. To be specific, if for example the length of one step is one year and r is the risk free interest rate for annual compounding, then X(n) = S(n)(1 + r)N −n − X(n − 1)(1 + r)N −n+1 . In order to introduce a general notion of derivative security which would cover futures, we first consider a natural generalisation X(n) = fn (S(0), S(1), . . . , S(n)) and then we push the level of generality ever further:

72


Definition 3.5 A derivative security (contingent claim) with the underlying asset represented by a sequence (S(n)) and the expiry time N is a sequence (X(1), . . . , X(N )) of random variables such that X(n) is measurable with respect to the σ-field Fn generated by (S(0), S(1), . . . , S(n)), for each n = 1, . . . , N.

Proposition 3.12 A derivative security X must be of the form X = f (S(0), S(1), . . . , S(N )) for some measurable f : RN +1 → R. We could make one more step and dispose of the underlying random variables. The role of the underlying object would be played by an increasing sequence of σ-fields Fn and we would say that a contingent claim (avoiding here the other term) is a sequence of random variables X(n) such that X(n) is Fn -measurable, but there is little need for such a generality in practical applications. The only case where that formulation would be relevant is the situation where there are no numerical observations but only some flow of information modelled by events and σ-fields.

Example 3.4 Payoffs of exotic options depend on the whole paths of consecutive stock prices. For example, the payoff of a European lookback option with exercise time N is determined by f (x0 , x1 , . . . , xN ) = max{x0 , x1 , . . . , xN } − xN

Exercise 3.9 Find the function f for a down-and-out call (which is a European call except that is ceases to exist if the stock price at any time before the exercise date goes below the barrier L < S(0)).

Example 3.5 Consider an American put option in a binomial model. We shall see that it fits the above abstract scheme. Recall that American options can be exercised at any time before expiry and the payoff of a put exercised at time n is (K−S(n))+ written g(S(n)) for brevity, g(x) = (K − x)+ . This option offers to the holder cash flow of the same nature as the stock. The latter is determined by the stock


73

price and stock can be sold at any time, of course only once. The American option can be sold or exercised also only once. The value of this option will be denoted by P A (n) we shall show that it is a derivative security in the sense of Definition 3.5. We shall demonstrate that it is possible to write P A (n) = fn (S(n)) for some functions fn . Consider an option expiring at N = 2. Clearly f2 (x) = g(x) At time n = 1 the holder of the option can exercise or wait till n = 2. The value of waiting is the same as the value of European put issued at n = 1 with exercise time N = 2 (which, as is well known and will be seen in Section 7.4.3 in some detail) can be computed as the expectation with respect to some probability p of the discounted payoff). The value of the American put is the greater of the two so o n 1 f1 (x) = max g(x), [pf2 (xU ) + (1 − p)f2 (xD)] . 1+r The same argument gives n o 1 f0 (x) = max g(x), [pf1 (xU ) + (1 − p∗ )f1 (xD)] . 1+r In general, for an American option expiring at time N we have the following chain of recursive formulae: fN (x) = g(x), n fn−1 (x) = max g(x),

o 1 [pfn (xU ) + (1 − p)fn (xD)] . 1+r

3.6 Proofs of propositions Proof (of Proposition 3.4) (i) We have proved that if f is measurable then so are f + , f − . Conversely, note that f (x) = f + (x) − f − (x) so Theorem 3.2 gives the result. (ii) The function u 7→ |u| is continuous so Lemma 3.3 with F (u, v) = |u| gives measurability of |f | (an alternative is to use |f | = f + + f − ). To see that the converse is not true take a non-measurable set A and let f = 1A − 1Ac . It is non-measurable since {x : f (x) > 0} = A is non-measurable. But |f | = 1 is clearly measurable.

74


Proof (of Proposition 3.9) Since f ≤ ess supf and g ≤ ess supg a.e., by adding we have f +g ≤ ess supf + ess sup g a.e. So the number ess sup f + ess sup g belongs to the set {z : f + g ≤ z a.e.} hence the infimum of this set is smaller than this number.

Proof (of Proposition 3.11) First note that the σ-field generated by S(N ) is of the form F = {S(N )−1 (B) : B -Borel} since these sets form a σ-field and any other σ-field such that S(N ) is measurable with respect to it has to contain all inverse images of Borel sets. Next we proceed in three steps: 1) Suppose X = 1A for A ∈ F. Then A = S(N )−1 (B) for a Borel subset of R. Put f = 1B and clearly X = f ◦ S(N ). P P 2) If X is a step function, X = ci 1Ai then take f = ci 1Bi where Ai = S(N )−1 (Bi ). 3) In general, a measurable function X can be approximated by step funcP22n (see Proposition 4.10 for more details) tions Xn = k=0 2kn · 1Y −1 ([ kn , k+1 2 2n )) and we take f = lim sup fn , where fn corresponds to Yn as in step 2) and the sequence clearly converges on the range of S(N ).

Proof (of Proposition 3.12) 1) Suppose X = 1A for A ∈ F. Then A = (S(1), . . . , S(N ))−1 (B) for Borel B ⊂ RN , and f = 1B satisfies the claim. Steps 2) and 3) are the same as in the proof of the previous proposition.

4 Integral

The theory developed below deals with Lebesgue measure for the sake of simplicity. However, all we need (except for the section where we discuss the Riemann integration) is the property of m being a measure, i.e. a countably additive (extended-) real valued function µ defined on a σ-field F of subsets of a fixed set Ω. Therefore, the theory developed for the measure space (R, M, m) in the following sections can be extended virtually without change to an abstractly given measure space (Ω, F, µ). We encourage the reader to bear in mind the possibility of such a generalization. We will need it in the probability section at the end of the chapter, and in the following chapters.

4.1 Definition of the integral We are now able to resolve one of the problems we identified earlier: how to integrate functions like 1Q , which take only finitely many values, but where the sets on which these values are taken are not at all ‘like intervals’.

Definition 4.1 A non-negative function ϕ : R → R which takes only finitely many values, i.e. the range of ϕ is a finite set of distinct non-negative reals {a1 , a2 , . . . , an }, is a 75

76


simple function if all the sets Ai = ϕ−1 ({ai }) = {x : ϕ(x) = ai },

i = 1, 2, . . . , n,

are measurable sets. Note that the sets Ai ∈ M are pairwise disjoint and their union is R. Clearly we can write ϕ(x) =

n X

ai 1Ai (x)

i=1

so that (by Theorem 3.2) each simple function is measurable.

Definition 4.2 The (Lebesgue) integral over E ∈ M of the simple function ϕ is given by: Z

ϕ dm = E

n X i=1

ai m(Ai ∩ E).

(Note: Since we shall allow m(Ai ) = +∞, we use the convention 0 × ∞ = 0 here.)

Figure 4.1 Integral of a simple function

Example 4.1 Consider the simple function 1Q which takes the value 1 on Q and 0 on R \ Q. By the above definition we have Z 1Q dm = 1 × m(Q) + 0 × m(R \ Q) = 0 R

4. Integral

77

since Q is a null set. Recall that this function is not Riemann-integrable. Similarly, 1C has integral 0, where C is the Cantor set.

Exercise 4.1 Find the integral of ϕ over E where (a) ϕ(x) = Int(x), E = [0, 10] (b) ϕ(x) = Int(x2 ), E = [0, 2] (c) ϕ(x) = Int(sin x), E = [0, 2π] and Int denotes the integer part of a real number. (Note: many texts use the symbol [x] to denote Int(x). We prefer to use Int for increased clarity.) In order to extend the integral to more general functions, Henri Lebesgue (in 1902) adopted an apparently obvious, but subtle device: instead of partitioning the domain of a bounded function f into many small intervals, he partitioned its range into a finite number of small intervals of the form Ai = [ai−1 , ai ), and approximated the ‘area’ under the graph of f by the upper sum S(n) =

n X

ai m(f −1 (Ai ))

i=1

and the lower sum s(n) =

n X

ai−1 m(f −1 (Ai ))

i=1

respectively; then integrable functions had the property that the infimum of all upper sums equals the supremum of all lower sums – mirroring Riemann’s construction (see also Figure 3.1). A century of experience with the Lebesgue integral has led to many equivalent definitions, some of them technically (if not always conceptually) simpler. We shall follow a version which, while very similar to Lebesgue’s original construction, allows us to make full use of the measure theory developed already. First we stay with non-negative functions:

Definition 4.3 For any non-negative measurable function f and E ∈ M the integral is defined as Z f dm = sup Y (E, f ) E

R

E

f dm

78


where Y (E, f ) =

nZ

E

o ϕ dm : 0 ≤ ϕ ≤ f, ϕ is simple .

Note that the integral can be +∞, and is always non-negative. Clearly, the set Y (E, f ) is always of the form [0, x] or [0, x), where the value x = +∞ is allowed. If E = [a, b] we write the integral as Z b Z b f dm, f (x) dm(x), a

a

Rb

Rb or Reven as a f (x) dx, when no confusion is possible (and we set a f dm = R R a − b f dm if a > b). The notation f dm means R f dm. Clearly, if for some A ∈ M and a non-negative measurable function g we have g = 0 on Ac , then any non-negative simple function that lies below g must be zero on Ac . Applying this to g = f.1A we obtain the important identity Z Z f dm = f 1A dm. A

Exercise 4.2 Suppose that f : [0, 1] → R is defined by letting f (x) = 0 on the Cantor −k set and f (x) = k for all x in each R 1 interval of length 3 which has been removed from [0, 1]. Calculate 0 f dm. Hint Recall that

P∞

k=1

kxk−1 =

d dx (

P∞

k=0

xk ) =

1 (1−x)2

when |x| < 1.

If f is a simple function, we now have two definitions of the integral; thus for consistency you should check carefully that the above definitions coincide.

Proposition 4.1 For simple functions, Definitions 4.2 and 4.3 are equivalent. Furthermore, we can prove the following basic properties of integrals of simple functions:

Theorem 4.2 Let ϕ, ψ be simple functions. Then:

4. Integral

79

(i) if ϕ ≤ ψ then

R

E

ϕ dm ≤

R

E

ψ dm,

(ii) if A, B are disjoint sets in M, then Z Z Z ϕ dm = ϕ dm + ϕ dm, A∪B

A

B

(iii) for all constants a > 0 Z

aϕ dm = a E

Z

ϕ dm. E

Proof (i) Notice that Y (E, ϕ) ⊆ Y (E, ψ) (we use Definition 4.3). P (ii) Employing the properties of m we have (ϕ = c i 1Di ) Z X ϕ dm = ci m(Di ∩ (A ∪ B)) A∪B X = ci m(Di ∩ A) + m(Di ∩ B) X X = ci m(Di ∩ A) + ci (Di ∩ B) Z Z = ϕ dm + ϕ dm. A

P

B

P

(iii) If ϕ = ci 1Ai then aϕ = aci 1Ai and Z Z X X aϕ dm = aci m(E ∩ Ai ) = a ci m(E ∩ Ai ) = a ϕ dm E

E

as required.

Next we show that the properties of the integrals of simple functions extend to the integrals of non-negative measurable functions:

Theorem 4.3 Suppose f and g are non-negative measurable functions. (i) If A ∈ M , and f ≤ g on A, then Z Z f dm ≤ g dm. A

A

(ii) If B ⊆ A, A, B ∈ M, then Z

B

f dm ≤

Z

f dm. A

80


(iii) For a ≥ 0,

Z

(iv) If A is null then

af dm = a A

Z

Z

f dm. A

f dm = 0. A

(v) If A, B ∈ M, A ∩ B = Ø, then Z Z Z f dm = f dm + f dm. A∪B

A

B

Proof (i) Notice that Y (A, f ) ⊆ Y (A, g) (there is more room to squeeze simple functions under g than under f ) and the sup of a bigger set is larger. (ii) If ϕ is a simple function lying below f on B, then extending it by zero outside B we obtain a simple function which is below f on A. The integrals of these simple functions are the same so Y (B, f ) ⊆ Y (A, f ) and we conclude as in (i). (iii) The elements of the set Y (A, af ) are of the form a × x where x ∈ Y (A, f ) so the same relation holds between R their suprema. P c i 1 Ei , (iv) For any simple function ϕ, A ϕ dm = 0. To see this, take ϕ = say, then m(A ∩ Ei ) = 0 for each i, so Y (A, f ) =R{0}. (v) The elements of Y (A ∪RB, f ) are ofR the form A∪B ϕ dm so by Theorem 4.2 (ii) they are of the form A ϕ dm + B Rϕ dm. So Y (A R Y (A, f ) + R ∪ B, f ) = Y (B, f ) and taking suprema this yields A∪B f dm ≤ A f dm + B f dm. For the opposite inequality, suppose that the simple functions ϕ and ψ satisfy: ϕ ≤ f on A and ϕ = 0 off A, while ψ ≤ f on B and ψ = 0 off B. Since A ∩ B = Ø, we can construct a new simple function γ ≤ f by setting γ = ϕ on A, γ = ψ on B and γ = 0 outside A ∪ B. Then Z Z Z Z ϕ dm + ψ dm = γ dm + γ dm A B B ZA = γ dm ZA∪B ≤ f dm. A∪B

On the right we have an upper bound which remains valid for all simple functions that lie below suprema over ϕ and ψ separately R R f on A ∪ R B. Thus taking on the left gives A f dm + B f dm ≤ A∪B f dm.

4. Integral

81

Exercise 4.3 Prove the following Mean RValue Theorem for the integral: if a ≤ f (x) ≤ b for x ∈ A, then am(A) ≤ A f dm ≤ bm(A). We now confirm that null sets are precisely the ‘negligible sets’ for integration theory:

Theorem 4.4 Suppose f is a non-negative measurable function. Then f = 0 a.e. if and only R if R f dm = 0.

Proof First, note that if f = 0 a.e. and 0 ≤ ϕ ≤ f is a simple R function, then ϕ = 0 a.e. since neither f nor ϕ take negative values. Thus R ϕ dm = 0 for all such R ϕ and so R f dm = 0Ralso. Conversely, given R f dm = 0, let E = {x : f (x) > 0}. Our goal is to show that m(E) = 0. Put 1 En = f −1 ([ , ∞)) for n ≥ 1. n Clearly, {En } increase to E with E=

∞ [

En .

n=1

To show that m(E) = 0 it is sufficient to prove that m(En ) = 0 for all n. (See Theorem 2.13.) The function ϕ = n1 1En is simple and ϕ ≤ f by the definition of En . So Z Z 1 ϕ dm = m(En ) ≤ f dm = 0 n R R hence m(En ) = 0 for all n.

Using the results proved so far the following ‘a.e.’ version of the monotonicity of the integral is not difficult to prove:

Proposition 4.5 If f and g are measurable then f ≤ g a.e. implies

R

f dm ≤

R

g dm.

82


Hint Let A = {x : f (x) ≤ g(x)}, then B = Ac is null and f 1A ≤ g1A . Now use Theorems 4.3 and 4.4. Using Theorems 3.2 and 3.5 you should now provide a second proof of a result we already noted in Proposition 3.4 but repeat here for emphasis:

Proposition 4.6 The function f : R → R is measurable iff both f + and f − are measurable.

4.2 Monotone Convergence Theorems The crux of Lebesgue integration is its convergence theory. We can make a start on that by giving a famous result

Theorem 4.7 (Fatou’s Lemma) If {fn } is a sequence of non-negative measurable functions then Z Z lim inf fn dm ≥ lim inf fn dm. n→∞

E

E

n→∞

Proof Write f = lim inf fn n→∞

and recall that f = lim gn n→∞

where gn = inf k≥n fk (the sequence gn is non-decreasing). Let ϕ be a simple function, ϕ ≤ f . To show that Z Z f dm ≤ lim inf fn dm E

n→∞

E

it is sufficient to see that Z for any such ϕ.

E

ϕ dm ≤ lim inf n→∞

Z

fn dm E

4. Integral

83

The set where f = 0 is irrelevant since it does not contribute to we can assume, without loss of generality, that f > 0 on E. Put ϕ(x) − ε > 0 if ϕ(x) > 0 ϕ(x) = 0 if ϕ(x) = 0 or x ∈ /E

R

E

f dm so

where ε is sufficiently small to ensure ϕ ≥ 0. Now ϕ < f , gn % f so ‘eventually’ gn ≥ ϕ. We make the last statement more precise: put Ak = {x : gk (x) ≥ ϕ(x)} and we have ∞ [

Ak ⊆ Ak+1 , Next, Z

An ∩E

ϕ dm ≤ ≤ ≤

Z

Z

Z

for k ≥ n. Hence

gn dm (as gn dominates ϕ on Ak ) An ∩E

fk dm

fk dm (as E is the larger set) E

Z

ϕ dm = An ∩E

l X i=1

for k ≥ n (by the definition of gn )

An ∩E

An ∩E

ϕ dm ≤ lim inf k→∞

Now we let n → ∞: writing ϕ = Z

Ak = R.

k=1

Pl

i=1 ci 1Bi

Z

fk dm.

(4.1)

E

for some ci ≥ 0, Bi ∈ M, i ≤ l

ci m(An ∩ E ∩ Bi ) −→

l X i=1

ci m(E ∩ Bi ) =

Z

ϕ dm E

and the inequality (4.1) remains true in the limit: Z Z ϕ dm ≤ lim inf fk dm. k→∞

E

E

We are close – all we need is to replace ϕ by ϕ in the last relation. This will be done by letting ε → 0 but some care will be needed. Suppose that m({x : ϕ(x) > 0}) < ∞. Then Z Z ϕ dm = ϕ dm − εm({x : ϕ(x) > 0}) E

E

and we get the result by letting ε → 0.

84

R


The case m({x R : ϕ(x) > 0}) = ∞ has to be treated separately. Here ϕ dm = ∞, so E f dm = ∞. We have to show that E Z lim inf fk dm = ∞. k→∞

E

Let ci be the values of ϕ and let a = 12 min{ci } ({ci } is a finite set!). Similarly to above put Dn = {x : gn (x) > a} and

Z

since Dn % R. As before Z

Dn ∩E

for k ≥ n, so lim inf

R

E

Dn ∩E

gn dm ≤

Z

gn dm → ∞

Dn ∩E

fk dm ≤

Z

fk dm E

fk dm has to be infinite.

Example 4.2 R Let fn = 1[n,n+1] . Clearly fn dm = 1 for all n, lim inf fn = 0 (= lim fn ), so the above inequality may be strict and we have Z Z (lim fn ) dm 6= lim fn dm.

Exercise 4.4 Construct an example of a sequence of functions with the strict inequality as above, such that all fn are zero outside the interval [0, 1]. It is now easy to prove one of the two main convergence theorems.

Theorem 4.8 (Monotone Convergence Theorem) If {fn } is a sequence of non-negative measurable functions, and {fn (x) : n ≥ 1} increases monotonically to f (x) for each x, i.e. fn % f pointwise, then Z Z lim fn (x) dm = f dm. n→∞

E

E

4. Integral

85

Proof Since fn ≤ f ,

R

E

fn dm ≤

R

f dm and so Z Z lim sup fn dm ≤ f dm. E

n→∞

Fatou’s lemma gives

Z

E

E

E

f dm ≤ lim inf n→∞

Z

fn dm E

which together with the basic relation Z Z lim inf fn dm ≤ lim sup fn dm n→∞

gives

n→∞

E

E

Z

Z Z f dm = lim inf fn dm = lim sup fn dm n→∞ n→∞ E E E R R hence the sequence E fn dm converges to E f dm.

Corollary 4.9 Suppose {fn } and f are non-negative Rand measurable. If {fn } increases to f R almost everywhere, then we still have E fn dm % E f dm for all measurable E.

Proof Suppose that fn % f a.e. and A is the set where the convergence holds, so that Ac is null. We can define fn on A gn = 0 on Ac , f on A g= 0 on Ac . Then using E = [E ∩ Ac ] ∪ [E ∩ A] we get Z Z Z gn dm = fn dm + 0 dm c E ZE∩A ZE∩A = fn dm + fn dm E∩Ac ZE∩A = fn dm E

c

R R (since E ∩A is null) and similarly E gRdm = E f dm. R The convergence gn → g holds everywhere so by Theorem 4.8, E gn dm → E g dm.

86


To apply the monotone convergence theorem it is convenient to approximate non-negative measurable functions by increasing sequences of simple functions.

Proposition 4.10 For any non-negative measurable f there is a sequence sn of non-negative simple functions such that sn % f . Hint Put

2n

2 X k sn = · 1 −1 k k+1 . 2n f ([ 2n , 2n )) k=0

Figure 4.2 Approximation by simple functions

4.3 Integrable functions All the hard work is done: we can extend the integral very easily to general real functions, using the positive part f + = max(f, 0), and the negative part f − = max(−f, 0), of any measurable function f : R → R. We will not use the non-negative measurable function |f | alone: as we saw in Proposition 3.4, |f | can be measurable without f being measurable!

Definition 4.4 R R If E ∈ M and the measurable function f has both E f + dm and E f − dm finite, then we say that f is integrable, and define Z Z Z f dm = f + dm − f − dm. E

E

E

4. Integral

87

The set of all functions that are integrable over E is denoted by L1 (E). In what follows E will be fixed and we often simply write L1 for L1 (E).

Exercise 4.5 For which α, is f (x) = xα in L1 (E) where (a) E = (0, 1); (b) E = (1, ∞)? Note that f is integrable iff |f | is integrable, and that Z Z Z |f | dm = f + dm + f − dm. E

E

E

Thus the Lebesgue integral is an ‘absolute’ integral: we cannot ‘make’ a function integrable by cancellation of large positive and negative parts. This has the consequence that some functions which have improper Riemann integrals fail to be Lebesgue integrable (see Section. 4.5). The properties of the integral of non-negative functions extend to any, not necessarily non-negative, integrable functions.

Proposition 4.11 If f and g are integrable, f ≤ g, then

R

f dm ≤

Hint If f ≤ g, then f + ≤ g + but f − ≥ g − .

R

g dm.

Remark 4.1 We observe (following [12], 5.12) that many proofs of results concerning integrable functions follow a standard pattern, utilising linearity and monotone convergence properties. To prove that a ’linear’ result holds for all functions in a space such as L1 (E) we proceed in four steps:

(i) verify that the required property holds for indicator functions – this is usually so by definition, (ii) use linearity to extend the property to non-negative simple functions, (iii) then use Monotone Convergence to show that the property is shared by all non-negative measurable functions,

(iv) finally, extend to the whole class of functions by writing f = f + − f − and using linearity again.

88


The next result gives a good illustration of the technique. R We wish to show that the mapping f 7→ A f dm is linear. This fact is interesting on its own, but will also allow us to show that L1 is a vector space.

Theorem 4.12 For any integrable functions f , g their sum f + g is also integrable and Z Z Z (f + g) dm = f dm + g dm. E

E

E

Proof We apply the technique described in Remark 4.1. Step 1. Suppose first that f and g are non-negative simple functions. The P P result is a matter of routine calculation: let f = a i 1A i , g = bj 1Bj . The sum f + g is also a simple function which can be written in the form X f +g = (ai + bj )1Ai ∩Bj . i,j

Therefore Z X (f + g) dm = (ai + bj )m(Ai ∩ Bj ∩ E) E

i,j

=

XX i

=

j

X i

=

X i

=

X i

=

X i

=

Z

ai

ai m(Ai ∩ Bj ∩ E) +

X j

ai m(

m(Ai ∩ Bj ∩ E) +

[ j

(Ai ∩ Bj ∩ E)) +

ai m(Ai ∩

[ j

Bj ∩ E) +

ai m(Ai ∩ E) +

f dm + E

XX

Z

X j

j

i

X

bj

j

X

bj m(Ai ∩ Bj ∩ E)

X i

bj m(

j

X j

m(Ai ∩ Bj ∩ E)

[ i

(Ai ∩ Bj ∩ E))

bj m(Bj ∩

[ i

Ai ∩ E)

bj m(Bj ∩ E)

g dm E

where we have used the additivity of m and the facts that Ai cover R and the same is true for Bj .

4. Integral

89

Step 2. Now suppose that f , g are non-negative measurable (not necessarily simple) functions. By Proposition 4.10 we can find sequences sn , tn of simple functions such that sn % f and tn % g. Clearly sn +tn % f +g hence using the monotone convergence theorem and the additivity property for simple functions we obtain Z Z (f + g) dm = lim (sn + tn ) dm n→∞ E E Z Z = lim sn dm + lim tn dm n→∞ E n→∞ E Z Z = f dm + g dm. E

E

This, in particular, implies that the integral of f + g is finite if the integrals of f and g are finite. Step 3. Finally, let f , g be arbitrary integrable functions. Since Z Z |f + g| dm ≤ (|f | + |g|) dm, E

E

we can use Step 2 to deduce that the left-hand side is finite. We have f + g = (f + g)+ − (f + g)−

f + g = (f + − f − ) + (g + − g − ) so (f + g)+ − (f + g)− = f + − f − + g + − g − . We rearrange the equality to have only additions on both sides (f + g)+ + f − + g − = f + + g + + (f + g)− . We have non-negative functions on both sides, so by what we have proved so far Z Z Z Z Z Z + − − + + (f + g) dm + f dm + g dm = f dm + g dm + (f + g)− dm E

E

E

E

E

E

hence Z Z Z Z Z Z + − + − + (f +g) dm− (f +g) dm = f dm− f dm+ g dm− g − dm. E

E

E

E

E

E

By definition of the integral the last relation implies the claim of the theorem.

90


The following result is a routine application of monotone convergence:

Proposition 4.13 If f is integrable and c ∈ R, then Z Z (cf ) dm = c f dm. E

E

Hint Approximate f by a sequence of simple functions. We complete the proof that L1 is a vector space:

Theorem 4.14 For any measurable E, L1 (E) is a vector space.

Proof Let f, g ∈ L1 . To show that f +g ∈ L1 we have to prove that |f +g| is integrable: Z Z Z Z |f + g| dm ≤ (|f | + |g|) dm = |f | dm + |g| dm < ∞. E

E

E

E

Now let c be a constant: Z Z Z |cf | dm = |c| |f | dm = |c| |f | dm < ∞ E

E

E

so that cf ∈ L1 (E). We can now answer an important question on the extent to which the integral determines the integrand.

Theorem 4.15 R R If A f dm R ≤ A g dm R for all A ∈ M, then f ≤ g almost everywhere. In particular, if A f dm = A g dm for all A ∈ M, then f = g almost everywhere.

Proof By additivity R of the integral (and Proposition 4.12 below) it is sufficient to show that A h dm ≥ 0 for all A ∈ M implies h ≥ 0 (and then take h = g − f ).

4. Integral

91

S Write A = {x : h(x) < 0}; then A = An where An = {x : h(x) ≤ − n1 }. By monotonicity of the integral Z Z 1 1 h dm ≤ − dm = − m(An ), n n An An which is non-negative but this can only happen if m(An ) = 0. The sequence of sets An increases with n, hence m(A) = R 0, and so h(x) ≥ 0 almost everywhere. A similar argument shows that if A h dm ≤ 0 for all A, then h R ≤ 0 a.e. This implies the second claim of the theorem: put h = g − f and A h dm is both non-negative and non-positive, hence h ≥ 0 and h ≤ 0 a.e. thus h = 0 a.e. The next Proposition lists further important properties of integrable functions, whose straightforward proofs are typical applications of the results proved so far.

Proposition 4.16 (i) An integrable function is a.e. finite. (ii) For measurable f and A m(A) inf f ≤ A

(iii) |

R

f dm| ≤

R

|f | dm.

(iv) Assume that f ≥ 0 and

R

Z

A

f dm ≤ m(A) sup f. A

f dm = 0. Then f = 0 a.e.

The following theorem gives us the possibility of constructing many interesting measures, and is essential for the development of probability distributions.

Theorem 4.17 Let f ≥ 0. Then A 7→

R

A

f dm is a measure.

Proof Denote µ(A) =

R

A

f dm. The goal is to show [ X µ( Ei ) = µ(Ei ) i

92


for pairwise disjoint Ei . To this end consider the sequence gn = f 1Sni=1 Ei and note that gn % g, where g = f 1S∞ . Now i=1 Ei Z

Z

gn dm =

Z

g dm = µ(

i=1

Ei ),

i=1

f dm = Sn

∞ [

Ei

n Z X i=1

f dm = Ei

n X

µ(Ei )

i=1

and the monotone convergence theorem completes the proof.

4.4 The Dominated Convergence Theorem Many questions in analysis centre on conditions under which the order of two limit processes, applied to certain functions, can be interchanged. Since integration is a limit process applied to measurable functions, it is natural to ask under what conditions on a pointwise (or pointwise a.e.) convergent sequence (fn ), the limit of the integrals isRthe integralR of the pointwise limit function f , i.e. when can we state that lim fn dm = (lim fn ) dm? The monotone convergence theorem (Theorem 4.8) provided the answer that this conclusion is valid for monotone increasing sequences of non-negative measurable functions, though in that case, of course, the limits may equal +∞. The following example shows that for general sequences of integrable functions the conclusion will not hold without some further conditions:

Example 4.3 Let fn (x) = n1[0, n1 ] (x). Clearly fn (x) → 0 for all x but

R

fn (x) dx = 1.

The limit theorem which turns out to be the most useful in practice states that convergence holds for an a.e. convergent sequence which is dominated by an integrable function. Again Fatou’s lemma holds the key to the proof.

Theorem 4.18 (Dominated Convergence Theorem) Suppose E ∈ M. Let (fn ) be a sequence of measurable functions such that |fn | ≤ g a.e. on E for all n ≥ 1, where g is integrable over E. If f = lim n→∞ fn a.e. then f is integrable over E and Z Z lim fn (x) dm = f dm. n→∞

E

E

4. Integral

93

Proof Suppose for the moment that fn ≥ 0. Fatou’s lemma gives Z Z f dm ≤ lim inf fn dm. n→∞

E

E

It is therefore sufficient to show that Z Z lim sup fn dm ≤ f dm. n→∞

E

(4.2)

E

Fatou’s lemma applied to g − fn gives Z Z lim (g − fn ) ≤ lim inf (g − fn ) dm. E n→∞

n→∞

E

On the left we have Z

E

(g − f ) dm =

Z

E

g dm −

Z

f dm. E

On the right lim inf

Z

(g − fn ) dm Z Z = lim inf g dm − fn dm n→∞ E Z Z E = g dm − lim sup fn dm, n→∞

E

E

n→∞

E

where we have used the elementary fact that lim inf (−an ) = − lim sup an . n→∞

n→∞

Putting this together we get Z Z Z Z g dm − f dm ≤ g dm − lim sup fn dm. E

Finally, subtract

E

R

E

E

n→∞

E

g dm (which is finite) and multiply by −1 to arrive at (4.2).

Now consider a general, not necessarily non-negative sequence (fn ). Since by the hypothesis −g(x) ≤ fn (x) ≤ g(x) we have 0 ≤ fn (x) + g(x) ≤ 2g(x) and we can apply the result proved for non-negative functions to the sequence fn (x) + g(x) (the function 2g is of course integrable).

94


Example 4.4 Going back to the example preceding the theorem, fn = n1[0, n1 ] , we can see that an integrable g to dominate fn cannot be found. The least upper bound 1 is g(x) = supn fn (x), g(x) = k on ( k+1 , k1 ] so Z

g(x) dx =

∞ X

k=1

∞

X 1 1 1 )= = +∞. k( − k k+1 k+1 k=1

For a typical positive example consider fn (x) =

n sin x 1 + n2 x1/2

R for x ∈ (0, 1). Clearly fn (x) → 0. To conclude that limn fn dm = 0 we need an integrable dominating function. This is usually where some ingenuity is needed; however in the present example the most straightforward estimate will suffice: n sin x n n 1 1 1 + n2 x1/2 ≤ 1 + n2 x1/2 ≤ n2 x1/2 = nx1/2 ≤ x1/2 .

(To see from first principles that the dominating function g : x 7→ √1x is integrable over [0, 1] can be rather tedious – cf. the worked example in Chapter 1 √ for the Riemann integral of x 7→ x. However, we shall show shortly that the Lebesgue and Riemann integrals of a bounded function coincide if the latter exists, and hence we can apply the Fundamental Theorem of the Calculus to confirm the integrability of g.) The following facts will be useful later.

Proposition 4.19 Suppose f is integrable and define gn = f 1[−n,n] , hn = min(f, n) (both truncate f in some R way: the gn vanish R outside a bounded interval, the hn are bounded). Then |f − gn | dm → 0, |f − hn | dm → 0. Hint Use the dominated convergence theorem.

Exercise 4.6 Use the dominated convergence theorem to find Z ∞ lim fn (x) dx n→∞

1

4. Integral

95

where fn (x) =

√

x . 1 + nx3

Exercise 4.7 Investigate the convergence of Z ∞ a

2

2

n2 xe−n x dx 1 + x2

for a > 0, and for a = 0.

Exercise 4.8 Investigate the convergence of Z ∞ 0

1

(1 +

x n n)

√ dx. n x

We will need the following extension of Theorem 4.12:

Proposition 4.20 For a sequence of non-negative measurable functions fn we have Z X ∞ ∞ Z X fn dm = fn dm. n=1

Hint The sequence gk =

n=1

Pk

n=1

fn is increasing and converges to

P∞

n=1

fn .

We cannot yet conclude that the sum of the series on the right-hand side is P∞ a.e. finite, so n=1 fn need not be integrable. However:

Theorem 4.21 (Beppo–Levi) Suppose that

Then the series and

P∞

k=1

∞ Z X

k=1

|fk | dm is finite.

fk (x) converges for almost all x, its sum is integrable, Z X ∞ k=1

fk dm =

∞ Z X

k=1

fk dm.

96


Proof The function ϕ(x) = sition 4.20

P∞

k=1

|fk (x)| is non-negative, measurable, and by Propo-

Z

ϕ dm =

∞ Z X k=1

|fk | dm.

This is finite, so ϕ is integrable. Therefore ϕ is finite a.e. Hence the series P∞ P∞ fk (x) converges (since it k=1 |fk (x)| converges a.e. and so the series k=1 P∞ converges absolutely) for almost all x. Let f (x) = k=1 fk (x) (put f (x) = 0 for x for which the series diverges – the value we choose is irrelevant since the set of such x is null). For all partial sums we have |

n X k=1

fk (x)| ≤ ϕ(x)

so we can apply the dominated convergence theorem to find Z Z n X f dm = lim fk dm n→∞

= lim

n→∞

= lim

n→∞

=

∞ Z X

k=1

Z X n

fk dm

k=1

n Z X

fk dm

k=1

fk dm

k=1

as required.

Example 4.5 P∞

1 kxk−1 = (1−x) 2 we can use the Beppo–Levi theorem R 1 log x 2 to evaluate the integral 0 ( 1−x ) dx : first let fn (x) = nxn−1 (log x)2 for n ≥ 1, x ∈ (0, 1), so that fn ≥ 0, fn is continuous, hence measurable, and P∞ log x 2 n=1 fn (x) = ( 1−x ) = f (x) is finite for x ∈ (0, 1). By Beppo–Levi the R1 R1 R1 P sum is integrable and 0 f (x) dx = ∞ n=1 0 fn (x) dx. To calculate 0 fn (x) dx R 1 n−1 we first use integration by parts to obtain 0 x (log x)2 dx = n23 . Thus R1 P∞ 1 2 f (x) dx = 2 n=1 n2 = π3. 0

Recalling that

k=1

Exercise 4.9 The following are variations on the above theme:

4. Integral

(a) For which values of a ∈ R does the power series integrable function on [−1, 1]? R∞ 2 (b) Show that 0 exx−1 dx = π6 .

97

P

n≥0

na xn define an

4.5 Relation to the Riemann integral Our prime motivation for introducing the Lebesgue integral has been to provide a sound theoretical foundation for the twin concepts of measure and integral, and to serve as the model upon which an abstract theory of measure spaces can be built. Such a general theory has many applications, a principal one being the mathematical foundations of the theory of probability. At the same time, Lebesgue integration has greater scope and more flexibility in dealing with limit operations than does its Riemann counterpart. However, just as with the Riemann integral, the computation of specific integrals from first principles is laborious, and we have, as yet, no simple ‘recipes’ for handling particular functions. To link the theory with the convenient techniques of elementary calculus we therefore need to take two further steps: to prove the Fundamental Theorem of the Calculus as stated in Chapter 1 and to show that the Lebesgue and Riemann integrals coincide whenever the latter exists. In the process we shall find necessary and sufficient conditions for the existence of the Riemann integral. In fact, given Proposition 4.16 the proof of the Fundamental Theorem becomes a simple application of the intermediate value theorem for continuous functions, and is left to the reader:

Proposition 4.22 If f : [a,Rb] → R is continuous then f is integrable and the function F given by x F (x) = a f dm is differentiable for x ∈ (a, b), with derivative F 0 = f .

R Hint Note that if f ∈ L1 and A, B ∈ M are disjoint, then A∪B f dm = R R x+h f dm + B f dm. Thus show that we can write F (x + h) − F (x) = x f dm A for fixed [x, x + h] ⊂ (a, b). R

We turn to showing that Lebesgue’s theory extends that of Riemann:

98


Theorem 4.23 Let f : [a, b] 7→ R be bounded.

(i) f is Riemann-integrable if and only if f is a.e. continuous with respect to Lebesgue measure on [a, b]. (ii) Riemann integrable functions on [a, b] are integrable with respect to Lebesgue measure on [a, b] and the integrals are the same.

Proof We need to prepare a little for the proof by recalling notation and some basic facts. Recall from Chapter 1 that any partition P = {ai : a = a0 < a1 < ... < an = b} of the interval [a, b], with ∆i = ai − ai−1 (i = 1, 2, . . . , n) and with Mi (resp. mi ) the sup (resp. inf) of f on Ii = [ai−1 , ai ], induces upper and lower Riemann Pn Pn sums UP = i=1 Mi ∆i and LP = i=1 mi ∆i . But these are just the Lebesgue Pn Pn integrals of the simple functions uP = i=1 Mi 1Ii and lP = i=1 mi 1Ii , by definition of the integral for such functions. Choose a sequence of partitions (Pn ) such that each Pn+1 refines Pn and the length of the largest subinterval in Pn goes to 0; writing un for uPn and ln for lPn we have ln ≤ f ≤ un for all n. Apply this on the measure space ([a, b], M[a,b] , m) where m = m[a,b] denotes Lebesgue measure restricted to [a, b]. Then u = inf n un and l = supn ln are measurable functions, and both sequences are monotone, since l1 ≤ l2 ≤ ... ≤ f ≤ ... ≤ u2 ≤ u1 .

(4.3)

Thus u = limn un and l = limn ln (pointwise) and all functions in (4.3) are bounded on [a, b] by M = sup{f (x) : x ∈ [a, b]}, which is integrable on [a, b]. By dominated convergence we conclude that Z b Z b Z b Z b lim Un = lim un dm = u dm, lim Ln = lim ln dm = l dm n

n

a

a

n

n

a

a

and the limit functions u and l are (Lebesgue-)integrable. Now suppose that x is not an endpoint of any of the intervals in the partitions (Pn ) – which excludes only countably many points of [a, b]. Then we have: f is continuous at x iff u(x) = f (x) = l(x). This follows at once from the definition of continuity, since the length of each subinterval approaches 0 and so the variation of f over the intervals containing x approaches 0 iff f is continuous at x.

4. Integral

99

Rb The Riemann integral a f (x) dx was defined as the common value of Rb Rb limn Un = a u dm and limn Ln = a ldm whenever these limits are equal. To prove (i), assume first Rthat f is Riemann-integrable, so that theRupper Rb b b and lower integrals coincide: a u dm = a l dm. But l ≤ f ≤ u, hence a (u − l) dm = 0 means that u = l = f a.e. by Theorem 4.15. Hence f is continuous a.e. by the above characterization of continuity of f at x, which only excludes a further null set of partition points. Conversely, if f is a.e. continuous, then u = f = l a.e. and u and l are Lebesgue-measurable, hence so is f (note that this uses the completeness of Lebesgue measure!). But f is also bounded by hypothesis, so it is Lebesgueintegrable over R[a, b], and as the integrals are a.e. equal, the integrals coincide b (but note that a f dm denotes the Lebesgue integral of f !): Z

b

l dm = a

Z

b

f dm = a

Z

b

u dm.

(4.4)

a

Since the outer integrals are the same, f is by definition also Riemannintegrable, which proves (i). To prove (ii), note simply that if f is Riemann-integrable, (i) shows that f is a.e. continuous, hence measurable, and then (4.4) shows that its Lebesgue integral coincides with the two outer integrals, hence with its Riemann integral.

Example 4.6 Recall the following example from Section 1.2: Dirichlet’s function defined on [0, 1] by 1 if x = m n n ∈Q f (x) = 0 if x ∈ /Q is a.e. continuous, hence Riemann-integrable, and its Riemann integral equals its Lebesgue integral, which is 0, since f is zero outside the null set Q. We have now justified the unproven claims made in earlier examples when evaluating integrals, since, at least for any continuous functions on bounded intervals, the techniques of elementary calculus also give the Lebesgue integrals of the functions concerned. Since the integral is additive over disjoint domains use of these techniques also extends to piecewise continuous functions.

Example 4.7 (Improper Riemann Integrals) Dealing with improper Riemann integrals involves an additional limit opera-

100


tion; we define such an integral by: Z ∞ f (x) dx := −∞

lim

a→−∞,b→∞

Z

b

f (x) dx a

whenever the double limit exists. (Other cases of ‘improper integrals’ are discussed in Remark 4.2.) Now suppose for the function f : R 7→ R this improper Riemann integral Rb exists. Then the Riemann integral a f (x) dx exists for each bounded interval [a, b], so that f is a.e. continuous on each [a, b], and thus on R. The converse is false, however: the function f which takes the value 1 on [n, n + 1) when n is even, and −1 when n is odd, is a.e. continuous (and thus Lebesgue measurable on R) but clearly the above limits fail to exist. More generally, it is not hard to show that if f ∈ L1 (R) then the above double limits will always exist. On the other hand, the existence of the double limit does not by itself guarantee that f ∈ L1 without further conditions: consider ( (−1)n if x ∈ [n, n + 1), n ≥ 0 n+1 f (x) = 0 if x < 0.

Figure 4.3 Graph of f Clearly the improper Riemann integral exists, Z ∞ ∞ X (−1)n f (x) dx = n+1 −∞ n=0 R P and the series converges. However, f ∈ / L1 , since R |f | dm = ∞ n=0 diverges.

1 n+1 ,

which

This yields another illustration of the ‘absolute’ nature of the Lebesgue integral: f ∈ L1 iff |f | ∈ L1 , so we cannot expect a finite sum for an integral whose ‘pieces’ make up a conditionally convergent series. For non-negative functions these problems do not arise; we have:

4. Integral

101

Theorem 4.24 If f ≥ 0 and R the above improper Riemann integral of f exists, then the Lebesgue integral R f dm always exists and equals the improper integral.

Proof To see this, simply note that the sequence (fn ) with fn = f 1[−n,n] increases monotonically to f , hence f is Lebesgue-measurable. Since fn is Riemannintegrable on [−n, n], the integrals coincide there, i.e. Z Z n fn dm = f (x) dx −n

R

1

for each n, so that fn ∈ L (R) for all n. By hypothesis the double limit Z n Z ∞ lim f (x) dx = f (x) dx n

−n

−∞

exists. On the other hand lim n

Z

fn dm = R

Z

f dm R

by monotone convergence, and so f ∈ L1 (R) and Z Z ∞ f dm = f (x) dx R

−∞

as required.

Exercise 4.10 Show that the function f given by f (x) = sinx x (x 6= 0) has an improper Riemann integral over R, but is not in L1 .

Remark 4.2 A second kind of improper Riemann integral is designed to handle functions which have asymptotes on a bounded interval, such as f (x) = x1 on (0, 1). For such cases we can define Z b Z b f (x) dx = lim f (x) dx a

ε&0

a+ε

when the limit exists. (Similar remarks apply to the upper limit of integration.)

102


4.6 Approximation of measurable functions The previous section provided an indication of the extent of the additional ‘freedom’ gained by developing the Lebesgue integral: Riemann integration binds us to functions whose discontinuities form an m-null set, while we can still find the Lebesgue integral of functions that are nowhere continuous, such as 1Q . We may ask, however, how real this additional generality is: can we, for example, approximate an arbitrary f ∈ L1 by continuous functions? In fact, since continuity is a local property, can we do this for arbitrary measurable functions? And this, in turn, provides a link with simple functions, since every measurable function is a limit of simple functions. We can go further, and ask whether for a simple function g approximating a given measurable function f we can choose the inverse image g −1 ({ai }) of each element of the range of g to P be an interval (such a g is usually called a step function; g = n cn 1In , where In are intervals). We shall tackle this question first:

Theorem 4.25 If f is a bounded measurable function on [a, b] and ε > 0 is given, then there Rb exists a step function h such that a |f − h| dm < ε.

Proof First assume additionally that f ≥ 0. Then sup{

Z

b a

Rb a

f dm is well-defined as

ϕ dm : 0 ≤ ϕ ≤ f, simple}.

Since f ≥ ϕ we have |f − ϕ| = f − ϕ, so we can find a simple function ϕ satisfying Z b Z b Z b ε |f − ϕ| dm = f dm − ϕ dm < . 2 a a a

It then remains to approximate an arbitrary simple function ϕ which vanishes off [a, b] by a step function h. The finite range {a1 , a2 , ..., an } of the function ϕ partitions [a, b], yielding disjoint measurable sets Ei = ϕ−1 ({ai }) such that Sn i=1 Ei = [a, b]. We now approximate each Ei by intervals: note that since ϕ is simple, M = sup{ϕ(x) : x ∈ [a, b]} < ∞. By Theorem 2.12 we can find open ε sets Oi such that Ei ⊂ Oi and m(Oi \ Ei ) < 2nM for i ≤ n. Since each Ei has finite measure, so do the Oi , hence each Oi can in turn be approximated by a S∞ finite union of disjoint open intervals: we know that Oi = j=1 Iij , where the P∞ open intervals can be chosen disjoint, so that m(Oi ) = j=1 m(Iij ) < ∞. As

4. Integral

103

Sk i ε the series converges, we can find ki such that m(Oi )−m( j=1 Iij ) < 2nM . Thus Rb Sk i ε with Gi = j=1 Iij we have a |1Ei − 1Gi | dm = m(Ei ∆Gi ) < nM for each Rb Pn i ≤ n. So set h = i=1 ai 1Gi . This step function satisfies a |ϕ − h| dm < 2ε Rb and hence a |f − h| dm < ε. The extension to general f is clear: f + and f − can be approximated to within ε2 by step functions h1 and h2 say, so with h = h1 − h2 we obtain Z

b a

|f − h| dm ≤

Z

b a

|f + − h1 | dm +

Z

b a

|f − − h2 | dm < ε

which completes the proof.

Figure 4.4 Approximation by continuous functions The ‘payoff’ is now immediate: with f and h as above, we can reorder the intervals Iij into a single finite sequence (Jm )m≤n with Jm = (cm , dm ) and h = Pn ε0 m=1 am 1Jm . We may assume that l(Jm ) = (dm − cm ) > 2 , and approximate 1Jm by a continuous function gm by setting gm = 1 on the slightly smaller 0 0 interval (cm + ε4 , d − ε4 ) and 0 outside Jm , while extending linearly in between Rb 0 (see Figure 4.4). It is obvious that gm is continuous and a |1Jm − gm | dm < ε2 . ε , where K = maxm≤n |am |, shows Repeating for each Jm and taking ε0 < nK Rb Pn ε that the continuous function g = m=1 am gm satisfies a |h − g| dm < 2 . Combining this inequality with Theorem 4.25 yields:

Theorem 4.26 Given f ∈ L1 and ε > 0, we canR find a continuous function g, vanishing outside some finite interval, such that |f − g| dm < ε.

104


Proof The preceding argument has verified this when f is a bounded measurable function vanishing off some interval [a, b]. For a given f ∈ L1 [a, b] we can again assume without loss that f ≥ 0. Let fn = min(f, n); Then R b the fn are bounded measurable functions dominated by f , fn → f , so that a |f − fN | dm < 2ε for some N . We R b can now find a continuous R b g, vanishing outside a finite interval, such that a |fN − g| dm < 2ε . Thus a |f − g|dm < ε. Finally, let f ∈ L1 (R) and f ≥ 0 be given. R R Choose n large enough to ensure that {|x|≥n} f dm < 3ε (which we can do as R |f | dm is finite; Proposition 4.19), R and simultaneously choose a continuous g with {|x|≥n} g dm < 3ε which satisfies Rn R ε −n |f − g| dm < 3 . Thus R |f − g| dm < ε. The well-known Riemann-Lebesgue lemma, which is very useful in the discussion of Fourier series, is easily deduced from the above approximation theorems:

Lemma 4.27 (Riemann-Lebesgue) R∞ 1 Suppose f ∈ L (R). Then the sequences s = k −∞ f (x) sin kx dx and ck = R∞ −∞ f (x) cos kx dx both converge to 0 as k → ∞.

Proof We prove this for (sk ) leaving the other, similar, case to the reader. For simR R∞ plicity of notation write for −∞ . The transformation x = y + πk shows that Z Z π π sk = f (y + ) sin(ky + π) dy = − f (y + ) sin(ky) dy. k k Since | sin x| ≤ 1, Z Z π π |f (x) − f (x + )| dx ≥ | (f (x) − f (x + )) sin kx dx| = 2|sk |. k k R It will therefore suffice to prove that |f (x) − f (x + h)| dx → 0 when h → 0. This is most easily done by approximating f by aRcontinuous g which vanishes outside some finite interval [a, b], and such that |f − g| dm < 3ε for a given ε > 0. For |h| < 1, the continuous function gh (x) = g(x + h) then vanishes off [a − 1, b + 1] and Z Z |f (x + h) − f (x)| dm ≤ |f (x + h) − g(x + h)| dm

4. Integral

105

+

Z

|g(x + h) − g(x)| dm +

Z

|g(x) − f (x)| dm.

The first and last integrals on the right are less than 3ε , while the integrand of ε whenever |h| < δ, by an appropriate the second can be made less than 3(b−a+2) choice of δ > 0, as g is continuous. As g vanishesRoutside [a − 1, b + 1], the second integral is also lessR than 3ε . Thus if |h| < δ, |f (x + h) − f (x) dm < ε. This proves that limk→∞ f (x) sin kx dx = 0.

4.7 Probability 4.7.1 Integration with respect to probability distributions Let X be a random variable with probability distribution PX . The following theorem shows how to perform a change of variable when integrating a function of X. In other words, it shows how to change the measure in an integral. This is fundamental in applying integration theory to probabilities. We emphasize again that only the closure properties of σ-fields and the countable additivity of measures are needed for the theorems we shall apply here, so that we can use an abstract formulation of a probability space (Ω, F, P ) in discussing their applications.

Theorem 4.28 Given a random variable X : Ω → R, Z Z g(X(ω)) dP (ω) = g(x) dPX (x). Ω

(4.5)

R

Proof We employ the technique described in Remark 4.1. For the indicator function g = 1A we have P (X ∈ A) on both sides. Then by linearity we have the result for simple functions. Approximation of non-negative measurable g by a monotone sequence of simple functions combined with the monotone convergence theorem gives the equality for such g. The case of general g ∈ L1 follows as before from the linearity of the integral, using g = g + − g − . The formula is useful in the case where the form of PX is known and allows one to carry out explicit computations.

106


Before we proceed to these situations, consider a very simple case as an illustration of the formula. Suppose that X is constant, i.e. X(ω) ≡ a. Then on the left in (4.5) we have the integral of a constant function, which equals g(a)P (Ω) = g(a) according to the general scheme of integrating indicator functions. On the right PX = δa and thusR we have a method of computing an integral with respect to Dirac measure: g(x) dδa = g(a). For discrete X taking values ai with probabilities pi we have Z X g(X) dP = g(ai )pi i

which is a well-known formula from elementary probability theory (see also P Section 3.5.3). In this case we have PX = i pi δai and on the right, the integral with respect to the combination of measures is the combination of the integrals: Z X Z g(x) dPX = pi g(x) dδai (x). i

In fact, this is a general property.

Theorem 4.29 If PX =

P

i

P pi Pi , where the Pi are probability measures, pi = 1, pi ≥ 0, then Z X Z g(x) dPX (x) = pi g(x) dPi . i

Proof The method is the same as above: first consider indicator functions 1A and the claim is just the definition of PX : on the left we have PX (A), on the right P i pi Pi (A). Then by additivity we get the formula for simple functions, and finally, approximation and use of the monotone convergence theorem completes the proof as before.

4.7.2 Absolutely continuous measures: examples of densities The measures P of the form A 7→ P (A) =

Z

f dm A

4. Integral

107

with non-negative integrable f will be called absolutely continuous, and the function f will be called a density of P with respect to Lebesgue measure, or simply a density. Clearly, for P to be a probability we have to impose the condition Z f dm = 1.

Students of probability often have an oversimplified mental picture of the world of random variables, believing that a random variable is either discrete or absolutely continuous. This image stems from the practical computational approach of many elementary textbooks, which present probability without the necessary background in measure theory. We have already provided a simple example which shows this to be a false dichotomy (Example 3.1). The simplest example of a density is this: let Ω ⊂ R be a Borel set with finite Lebesgue measure and put ( 1 if x ∈ Ω m(Ω) f (x) = 0 otherwise. We have already come across this sort of measure in the previous chapter, that is, the probability distribution of a specific random variable. We say that in this case the measure (distribution) is uniform. It corresponds to the case where the values of the random variable are spread evenly across some set, typically an interval, such as in choosing a number at random (Example 2.2). Slightly more complicated is the so-called triangle distribution with the density of the form shown in Figure 4.5.

Figure 4.5 Triangle distribution The most famous is the Gaussian or normal density n(x) = √

(x−µ)2 1 e− 2σ2 . 2πσ

(4.6)

This function is symmetric with respect to x = µ, and vanishes at infinity, i.e. limx→−∞ n(x) = 0 = limx→∞ n(x).

108


Figure 4.6 Gaussian distribution

Exercise 4.11 Show that

R∞

−∞

n(x) dx = 1.

Hint First consider the case µ = 0, σ = 1 and then transform the general case to this. The meaning of the number µ will become clear below and σ will be explained in the next chapter. Another widely used example is the Cauchy density: c(x) =

1 1 . π 1 + x2

This density gives rise to many counterexamples to ‘theorems’ which are too good to be true.

Exercise 4.12 Show that

R∞

−∞

c(x) dx = 1.

The exponential density is given by −λx ce f (x) = 0

if x ≥ 0 otherwise.

Exercise 4.13 Find the constant c for f to be a density of probability distribution.

4. Integral

109

The gamma distribution is really a large family of distributions, indexed by a parameter t > 0. It contains the exponential distribution as the special case where t = 1. Its density is defined as ( 1 t t−1 −λx e if x ≥ 0 Γ (t) λ x f (x) = 0 otherwise R∞ where the gamma function Γ (t) = 0 xt−1 e−x dx. The gamma distribution contains another widely used distribution as a special case: the distribution obtained from the density f when λ = 21 and t = d2 for some d ∈ N is denoted by χ2 (d) and called the chi-squared distribution with d degrees of freedom. The (cumulative) distribution function corresponding to a density is given by F (y) =

Z

y

f (x) dx. −∞

If f is continuous then F is differentiable and F 0 (x) = f (x) by the Fundamental Theorem of Calculus (see Proposition 4.22). We say that F is absolutely continuous if this relation holds with integrable f , and then f is the density of the probability measure induced by F . The following example due to Lebesgue shows that continuity of F is not sufficient for the existence of a density.

Example 4.8 Recall the Lebesgue function F define on page 20. We have F (y) = 0 for y ≤ 0, F (y) = 1 for y ≥ 1, F (y) = 12 for y ∈ [ 13 , 23 ), F (y) = 41 for y ∈ [ 19 , 92 ), F (y) = 43 for y ∈ [ 79 , 89 ) and so on. The function F is constant on the intervals removed in the process of constructing the Cantor set.

Figure 4.7 Lebesgue’s function

110


It is differentiable almost everywhere and the derivative is zero. So F cannot be absolutely continuous since then f would be zero almost everywhere, but on the other hand its integral is 1. We now define the (cumulative) distribution function of a random variable X : Ω → R, where, as above, (Ω, F, P ) is a given probability space: FX (y) = P ({ω : X(ω) ≤ y}) = PX ((−∞, y]).

Proposition 4.30 (i) FX is non-decreasing (y1 ≤ y2 implies FX (y1 ) ≤ FX (y2 )),

(ii) limy→∞ FX (y) = 1, limy→−∞ FX (y) = 0,

(iii) FX is right continuous (if y → y0 , y ≥ y0 , then FX (y) → F (y0 )).

Exercise 4.14 Show that FX is continuous if and only if PX ({y}) = 0 for all y.

Exercise 4.15 Find FX for (a) a constant random variable X, X(ω) = a for all ω (b) X : [0, 1] → R given by X(ω) = min{ω, 1 − ω} (the distance to the nearest endpoint of the interval [0, 1]) (c) X : [0, 1]2 → R, the distance to the nearest edge of the square [0, 1]2 . The fact that we are doing probability on subsets of Rn as sample spaces turns out to be not restrictive. In fact, the interval [0, 1] is sufficient as the following Skorokhod representation theorem shows.

Theorem 4.31 If a function F : R → [0, 1] satisfies conditions (i)–(iii) of Proposition 4.30, then there is a random variable defined on the probability space ([0, 1], B, m[0,1]), X : [0, 1] → R, such that F = FX .

4. Integral

111

Proof We write, for ω ∈ [0, 1], X + (ω) = inf{x : F (x) > ω},

X − (ω) = sup{x : F (x) < ω}.

Figure 4.8 Construction of X − ; continuity point

Figure 4.9 Construction of X − ; discontinuity point Three possible cases are illustrated in Figures 4.8, 4.9 and 4.10. We show that FX − = F , and for that we have to show that F (y) = m({ω : X − (ω) ≤ y}). The set {ω : X − (ω) ≤ y} is an interval with left endpoint 0. We are done if we show that its right endpoint is F (y), i.e. if X − (ω) ≤ y is equivalent to ω ≤ F (y). Suppose that ω ≤ F (y). Then {x : F (x) < ω} ⊂ {x : F (x) < F (y)} ⊂ {x : x ≤ y} (the last inclusion by the monotonicity of F ), hence X − (ω) = sup{x : F (x) < ω} ≤ y. Suppose that X − (ω) ≤ y. By monotonicity F (X − (ω)) ≤ F (y). By the rightcontinuity of F , ω ≤ F (X − (ω)) (if ω > F (X − (ω)), then there is x0 > X − (ω)

112


Figure 4.10 Construction of X − ; ‘flat’ piece such that F (X − (ω)) < F (x0 ) < ω, which is impossible since x0 is in the set whose supremum is taken to get X − (ω)) so ω ≤ F (y). For future use we also show that FX + = F . It is sufficient to see that m({ω : X − (ω) < X + (ω)}) = 0 (which is intuitively clear as this may happen only when the graph of F is ‘flat’, and there are countably many values corresponding to the ‘flat’ pieces, their Lebesgue measure being zero). More rigorously, [ {ω : X − (ω) < X + (ω)} = {ω : X − (ω) ≤ q < X + (ω)} q∈Q

and m({ω : X − (ω) ≤ q < X + (ω)}) = m({ω : X − (ω) ≤ q} \ {ω : X + (ω) ≤ q}) = F (q) − F (q) = 0. The following theorem provides a powerful method for calculating integrals relative to absolutely continuous distributions. The result holds for general measures but we formulate it for a probability distribution of a random variable in order not to overload or confuse the notation.

Theorem 4.32 If PX defined on Rn is absolutely continuous with density fX , g : Rn → R is integrable with respect to PX , then Z Z g(x) dPX (x) = fX (x)g(x) dx. Rn

Rn

Proof For an indicator function g(x) = 1A (x) we have PX (A) on the left which equals

4. Integral

R

113

R by the form of P , and consequently is equal to Rn 1A (x)fX (x) dx, i.e. the right-hand side. Extension to simple functions by linearity and to general integrable g by limit passage is routine. A fX (x) dx

Corollary 4.33 In the situation of the previous theorem we have Z Z g(X) dP = fX (x)g(x) dx. Ω

Rn

Proof This is an immediate consequence of the above theorem and Theorem 4.28. We conclude this section with a formula for a density of a function of a random variable with given density. Suppose that fX is known and we want to find the density of Y = g(X).

Theorem 4.34 If g : R → R is increasing and differentiable (thus invertible), then fg(X) (y) = fX (g −1 (y))

d −1 g (y). dy

Proof Consider the distribution function: Fg(X) (y) = P (g(X) ≤ y) = P (X ≤ g −1 (y)) = FX (g −1 (y)). Differentiate with respect to y to get the result.

Remark 4.3 A similar result holds if g is decreasing. The same argument as above gives fg(X) (y) = −fX (g −1 (y))

d −1 g (y). dy

114


Example 4.9 If X has standard normal distribution 1 1 2 n(x) = √ e− 2 x 2π (i.e. µ = 0 and σ = 1 in (4.6)), then the density of Y = µ + σX is given by (4.6). This follows at once from Theorem 4.34: g −1 (y) = µ−x σ ; its derivative is equal to σ1 .

Exercise 4.16 Find the density of Y = X 3 where fX = 1[0,1] .

4.7.3 Expectation of a random variable If X is a random variable defined on a probability space (Ω, F, P ) then we introduce the following notation: Z E(X) = X dP Ω

and we call this abstract integral the mathematical expectation of X. Using the results from the previous section we immediately have the following formulae: the expectation can be computed using the probability distribution: Z ∞ E(X) = x dPX (x), −∞

and for absolutely continuous X we have Z ∞ E(X) = xfX (x) dx. −∞

Example 4.10 Suppose that PX = 12 P1 + 12 P2 , where P1 = δa , P2 has a density f2 . Then Z 1 1 E(X) = a + xf (x) dx. 2 2

So, going back to Example 3.1 we can compute the expectation of the random variable considered there: Z 25 1 1 1 E(X) = · 0 + x dx = 6.25. 2 2 25 0

4. Integral

115

Exercise 4.17 Find the expectation of (a) a constant random variable X, X(ω) = a for all ω (b) X : [0, 1] → R given by X(ω) = min{ω, 1 − ω} (the distance to the nearest endpoint of the interval [0, 1]) (c) X : [0, 1]2 → R, the distance to the nearest edge of the square [0, 1]2 .

Exercise 4.18 Find the mathematical expectation of a random variable with (a) uniform distribution over the interval [a, b], (b) triangle distribution, (c) exponential distribution.

4.7.4 Characteristic function In what follows we will need the integrals of some complex functions. The theory is a straightforward extension of the real case. Let Z = X + iY where X, Y are real-valued random variables and define Z Z Z Z dP = X dP + i Y dP. Clearly, linearity of the integral and the dominated convergence theorem hold for the complex case. Another important relation which remains true is: Z Z | Z dP | ≤ |Z| dP. R R To see this consider the polar decomposition of Z dPR= | Z dP |e−iθ R . Then, iθ with 0. Then the simple functions sn = n1 R A satisfy sn ≤ f , but sn dm = nm(A) and the supremum here is ∞. Thus f dm = ∞ – a contradiction.

4. Integral

121

(ii) The simple function s(x) = c1A with c = inf A f has integral inf A f m(A) and satisfies s ≤ f , which R proves Rthe first inequality. Put t(x) = d1A with d = supA f and f ≤ t so f dm ≤ t dm which is the second inequality. (iii) Note that −|f | ≤ f ≤ |f | hence − are done.

R

|f | dm ≤

R

f dm ≤

R

|f | dm and we

S∞ (iv) Let En = f −1 ([ n1 , ∞)), and E = n=1 En . The sets Ei are measurable 1 and so is RE. The function function with s ≤ f . Hence R R s = n 1En is a simple s dm ≤ f dm = 0, so sn dm = 0, hence n1 m(En ) = 0. Finally, m(En ) = 0 for all n. Since En ⊂ En+1 , m(E) = lim m(En ) = 0. But E = {x : f (x) > 0} so f is zero outside the null set E.

Proof (of Proposition 4.19) If n → ∞ then 1[−n,n] → 1 hence gn = f 1[−n,n] → f . The convergence is R dominated: gn ≤ |f | and by the dominated convergence theorem we have R |f − gn | dm → 0. Similarly, hn = min(f, n) → f as n → ∞ and hn ≤ |f | so |f − hn | dm → 0.

Proof (of Proposition 4.20) Using

R

(f + g) dm =

R

f dm + Z X n k=1

R

g dm we can easily obtain (by induction)

fk dm =

n Z X

fk dm

k=1

Pn for any n. The sequence k=1 fk is increasing (fk ≥ 0) and converges to P∞ f . So the monotone convergence theorem gives k=1 k Z X ∞

fk dm = lim

k=1

n→∞

Z X n k=1

fk dm = lim

n→∞

n Z X

k=1

fk dm =

∞ Z X

fk dm

k=1

as required.

Proof (of Proposition 4.22) Continuous functions are measurable, and f is bounded on [a, b], hence f ∈ R x+h L1 [a, b]. Fix a < x < x + h < b, then F (x + h) − F (x) = x f dm, since the intervals [a, x] and (x, x + h] are disjoint, so that the integral is additive with respect to the upper endpoint. By the mean value property the values of righthand integrals are contained in the interval [Ah, Bh], where A = inf{f (t) : t ∈

122


[x, x + h]} and B = sup{f (t) : t ∈ [x, x + h]}. Both extrema are attained, as f is continuous, so we can find t1 , t2 in [x, x + h] with A = f (t1 ), B = f (t2 ). Thus Z 1 x+h f (t1 ) ≤ f dm ≤ f (t2 ). h x

The value theorem provides θ ∈ [0, 1] such that f (x + θh) = R intermediate F (x+h)−F (x) 1 x+h f dm = . Letting h → 0, the continuity of f ensures that h x h F 0 (x) = f (x).

Proof (of Proposition 4.30) (i) If y1 ≤ y2 , then {ω : X(ω) ≤ y1 } ⊂ {ω : X(ω) ≤ y2 } and by the monotonicity of measure FX (y1 ) = P ({ω : X(ω) ≤ y1 }) ≤ P ({ω : X(ω) ≤ y2 }) = FX (y2 ). S (ii) Let n → ∞; then n {ω : X(ω) ≤ n} = Ω (the sets increase). Hence P ({ω : X(ω) ≤ n}) → P (Ω) = 1 by Theorem 2.13 (i) and so limy→∞ FX (y) = 1. For the second claim consider FX (−n) = P ({ω : X(ω) ≤ −n}) and note T that limy→−∞ FX (y) = P ( n {ω : X(ω) ≤ −n}) = P (Ø) = 0.

(iii) This follows directly from Theorem 2.13 (ii) with An = {ω : X(ω) ≤ yn }, T yn % y, because FX (y) = P ( n {ω : X(ω) ≤ yn }).

Proof (of Proposition 4.36) Inserting the formulae for the option prices in call-put parity we have Z Z + S(0) = exp{−rT } (S(T ) − K) dP − (K − S(T ))+ dP + K Ω Ω Z = exp{−rT } (S(T ) − K)dP {S(T )≥K} Z − (K − S(T ))dP + K {S(T ) 0 ∃N : ∀n, m ≥ N kfn − fm kX < ε. If each Cauchy sequence is convergent to some element of X, then we say that X is complete.

5. Spaces of integrable functions

129

Example 5.1 Let fn (x) = x1 1[n,n+1] (x), and suppose that n 6= m. Z ∞ 1 kfn − fm k1 = |1[n,n+1] − 1[m,m+1] | dx x 0 Z n+1 Z m+1 1 1 dx + dx = x x m n n+1 m+1 = log + log . n m If an → 1, then log an → 0 and the right-hand side can be as small as we wish: for ε > 0 take N such that log NN+1 < 2ε . So fn is a Cauchy sequence in L1 (0, ∞). (When E = (a, b), we write L1 (a, b) for L1 (E), etc.)

Exercise 5.2 Decide whether each of the following is Cauchy as a sequence in L1 (0, ∞) (a) fn = 1[n,n+1] (b) fn = x1 1(0,n) (c) fn =

1 x2 1(0,n)

The proof of the main result below makes essential use of the Beppo–Levi theorem in order to transfer the main convergence question to that of series of real numbers; its role is essentially to provide the analogue of the fact that in R (and hence in C) absolutely convergent series will always converge. (The Beppo–Levi theorem clearly extends to complex-valued functions, just as we showed for the dominated convergence theorem, but we shall concentrate on the real case in the proof below, since the extension to C is immediate.) We digress briefly to recall how this property of series ensures completeness in R: let (xn ) be a Cauchy sequence in R, and extract a subsequence (xnk ) such that |xn − xnk | < 2−k for all n ≥ nk as follows: 1. find n1 such that |xn − xn1 | < 2−1 for all n ≥ n1 ,

2. find n2 > n1 such that |xn − xn2 | < 2−2 for all n ≥ n2 ,

3. . . .

4. find nk > nk−1 with |xn − xnk | < 2−k for all n ≥ nk .

The Cauchy property ensures each time that such nk can be found. Now consider the telescoping series with partial sums yk = xn1 + (xn2 − xn1 ) + · · · + (xnk − xnk−1 ) = xnk

130


which has |yk | ≤ |xn1 | +

k X i=1

|xni − xni−1 | < |xn1 | +

k X 1 . 2i i=1

Thus this series converges, in other words (xnk ) converges in R, and its limit is also that of the whole Cauchy sequence (xn ). To apply the Beppo–Levi theorem below we therefore need to extract a ‘rapidly convergent sequence’ from the given Cauchy sequence in L1 (E). This provides an a.e. limit for the original sequence, and the Fatou lemma does the rest.

Theorem 5.1 The space L1 (E) is complete.

Proof Suppose that fn is a Cauchy sequence. Let ε = n ≥ N1 1 kfn − fN1 k1 ≤ . 2 Next, let ε =

1 22 ,

1 2.

There is N1 such that for

and for some N2 > N1 we have kfn − fN2 k1 ≤

1 22

for n ≥ N2 . In this way we construct a subsequence fNk satisfying kfNn+1 − fNn k1 ≤ for all n. Hence the series Levi theorem, the series

P

n≥1

fN1 (x) +

1 2n

kfNn+1 − fNn k1 converges and by the Beppo– ∞ X

n=1

[fNn+1 (x) − fNn (x)]

converges a.e.; denote the sum by f (x). Since fN1 (x) +

k X

n=1

[fNn+1 (x) − fNn (x)] = fNk+1

the left-hand side converges to f (x), so fNk+1 (x) converges to f (x). Since the sequence of real numbers fn (x) is Cauchy and the above subsequence converges, the whole sequence converges to the same limit f (x).


131

We have to show that f ∈ L1 and kfk − f k1 → 0. Let ε > 0. The Cauchy condition gives an N such that ∀n, m ≥ N, kfn − fm k1 < ε. By Fatou’s lemma Z Z kf −fm k1 = |f −fm | dm ≤ lim inf |fNk −fm | dm = lim inf kfNk −fm k1 < ε. k→∞

k→∞

(5.1) So f − fm ∈ L1 which implies f = (f − fm ) + fm ∈ L1 , but (5.1) also gives kf − fm k1 → 0.

5.2 The Hilbert space L2 The space we now introduce plays a special role in the theory. It provides the closest analogue of the Euclidean space Rn among the spaces of functions, and its geometry is closely modelled on that of Rn . It is possible, via the integral, to induce the norm via an inner product, which in turn provides a concept of orthogonality (and hence ‘angles’) between functions. This gives L2 many pleasant properties, such as a ‘Pythagoras theorem’ and the concept of orthogonal projections, which plays vital role in many applications. To define the norm, and hence the space L2 (E) for a given measurable set E ⊂ R, let Z kf k2 = (

2

1

E

|f |2 dm) 2

and define L (E) as the set of measurable functions for which this quantity is finite. (Note that, as for L1 , we require non-negative integrands; it is essential that the integral is non-negative in order for the square root to make sense. Although we always have f 2 (x) = (f (x))2 ≥ 0 when f (x) is real, the modulus is needed to include the case of complex-valued functions f : E → C. This also makes the notation consistent with that of the other Lp -spaces we shall consider below where |f |2 is replaced by |f |p for arbitrary p ≥ 1.) We introduce L2 (E) as the set of equivalence classes of elements of L2 (E), under the equivalence relation f ≡ g iff f = g a.e., exactly as for L1 (E), and continue the convention of treating the equivalence classes as functions. If R f : E → C satisfies E |f |2 dm < ∞ we write f ∈ L2 (E, C) – again using f interchangeably as a representative of its equivalence class and to denote the class itself.

132


It is straightforward to prove that L2 (E) is a vector space: clearly, for a ∈ R, |af |2 is integrable if |f |2 is, while |f + g|2 ≤ 22 max{|f |2 , |g|2 } ≤ 4(|f |2 + |g|2 )

shows that L2 (E) is closed under addition.

5.2.1 Properties of the L2 -norm We provide a simple proof that the map f 7→ kf k2 is a norm: to see that it satisfies the triangle inequality requires a little work, but the ideas will be very familiar from elementary analysis in Rn , as is the terminology, though the context is rather different. We state and prove the result for the general case of L2 (E, C).

Theorem 5.2 (The Schwarz Inequality) If f, g ∈ L2 (E, C) then f g ∈ L1 (E, C) and Z | f g dm| ≤ kf gk1 ≤ kf k2 kgk2

(5.2)

E

where g denotes the complex conjugate of g.

Proof Replacing f , g by |f |, |g| we may assume that f and gR are non-negative (the first inequality has already been verified, since kf gk1 = E |f g| dm, and the second only Rinvolves the modulus in each case). Since we do not know in advance that E f g dm is finite, we shall first restrict attention to bounded measurable functions by setting fn = f ∧ n and gn = g ∧ n, and confine our domain of integration to the bounded set E ∩ [−k, k] = Ek . For any t ∈ R we have Z Z Z Z 2 2 2 0≤ (fn + tgn ) dm = fn dm + 2t fn gn dm + t gn2 dm. Ek

Ek

Ek

Ek

As a quadratic in t this does not have two distinct solutions, so the discriminant is non-positive. Thus for all n ≥ 1 Z Z Z (2 fn gn dm)2 ≤ 4( fn2 dm)( gn2 dm)2 Ek

≤ 4(

Z

Ek 2 E

|f |2 dm)(

= kf k22 kgk22 .

Ek 2

Z

E

|g|2 dm)


133

Monotone convergence now yields Z ( f g dm)2 ≤ kf k22kgk22 Ek

for each k, and since E =

S (

k

Z

Ek we obtain finally that

E

f g dm)2 ≤ kf k22kgk22 ,

which implies the Schwarz inequality. The triangle inequality for the norm on L2 (E, C) now follows at once – we need to show that kf + gk2 ≤ kf k2 + kgk2 for f, g ∈ L2 (E, C): Z Z Z kf + gk22 = |f + g|2 dm = (f + g)(f + g)dm = (f + g)(f + g) dm. E

E

E

The latter integral is Z Z Z |f |2 dm + (f g + f g) dm + |g|2 dm, E

E

E

which is dominated by (kf k2 + kgk2 )2 since the Schwarz inequality gives Z Z (f g + gf ) dm ≤ 2 |f g| dm ≤ 2kf k2kgk2 . E

E

The result follows. The other properties are immediate: (i) clearly kf k2 = 0 means that |f |2 = 0 a.e., hence f = 0 a.e., R2 1 (ii) for a ∈ C, kaf k2 = ( E |af |2 dm) 2 = |a|kf k2 .

Thus the map f 7→ kf k2 is a norm on L2 (E, C). The proof that L2 (E) is complete under this norm is similar to that for L1 (E), and will be given in Theorem 5.11 below for arbitrary Lp -spaces (1 < p < ∞).

In general, without restriction of the domain set E, neither L1 ⊆ L2 nor L ⊆ L1 . To see this consider E = [1, ∞), f (x) = x1 . Then f ∈ L2 (E) but f∈ / L1 (E). Next put F = (0, 1), g(x) = √1x . Now g ∈ L1 (F ) but g ∈ / L2 (F ). For finite measure spaces – and hence for probability spaces! – we do have a useful inclusion: 2

Proposition 5.3 If the set D has finite measure (that is, m(D) < ∞), then L2 (D) ⊂ L1 (D).

134


Hint Estimate |f | by means of |f |2 and then use the fact that the integral of |f |2 is finite. Before exploring the geometry induced on L2 by its norm, we consider examples of sequences in L2 to provide a little practice in determining which are Cauchy sequences for the L2 -norm, and compare this with their behaviour as elements of L1 .

Example 5.2 We show that the sequence fn = x1 1[n,n+1] is Cauchy in L2 (0, ∞). Z ∞ 1 |1 − 1[m,m+1] |2 dx kfn − fm k2 = 2 [n,n+1] x 0 Z n+1 Z m+1 1 1 = dx + dx 2 2 x x m n 1 1 1 1 = − + − n n+1 m m+1 2 2 ≤ + n m and for ε > 0 let N be such that m, n ≥ N .

2 N

kh00 k2 unless k = h0 . Hence kh00 k = kh − h0 k = inf{kh − kk : k ∈ K} = δK , say. Conversely, having found h0 ∈ K such that kh − h0 k = δK then for any real t and k ∈ K, h0 + tk ∈ K, so that kh − (h0 + tk)k2 ≥ kh − h0 k2 . Multiplying out the inner products and writing h00 = h − h0 , this means that −t[(h00 , k)+(k, h00 )]+t2 kkk2 ≥ 0. This can only hold for all t near 0 if (h00 , k) = 0, so that h00 ⊥k for every k ∈ K. To find h0 ∈ K with kh − h0 k = δK , first choose a sequence (kn ) in K such that kh − kn k → δK as n → ∞. Then apply the parallelogram law (Proposition 5.5 (i)) to the vectors h1 = h− 21 (km + kn ) and h2 = 12 (km − kn ). Note that h1 + h2 = h − kn and h1 − h2 = h − km . Hence the parallelogram law reads 1 1 kh − kn k2 + kh − km k2 = 2(kh − (km + kn )k2 + k (km − kn )k2 ) 2 2 2 and since 21 (km +kn ) ∈ K, kh− 12 (km +kn )k2 ≥ δK . As m, n → ∞ the left-hand 2 side converges to 2δK , hence that final term on the right must converge to 0. Thus the sequence (kn ) is Cauchy in K, and so converges to an element h0 of K. But since kh − kn k → δK while kkn − h0 k → 0 as n → ∞, kh − h0 k ≤ kh − kn k + kkn − h0 k shows that kh − h0 k = δK . This completes the proof.

In writing h = h0 + h00 we have decomposed the vector h ∈ H as the sum of two vectors, the first being its orthogonal projection onto K, while the second is orthogonal to all vectors in K. We say that h00 is orthogonal to K, and denote the set of all vectors orthogonal to K by K ⊥ . This exhibits H as a direct sum H = K ⊕ K ⊥ with each vector of the first factor being orthogonal to each vector in the second factor. We shall use the existence of orthogonal projections onto subspaces of L2 (Ω, F, P ) to construct the conditional expectation of a random variable with respect to a σ-field in Section 5.4.3.

Remark 5.2 The foregoing discussion barely scratches the surface of the structure of inner product spaces, such as L2 (E), which is elegantly explained, for example in [10]. On the one hand, the concept of orthogonality in an inner product space leads to consideration of orthonormal sets, i.e. families (eα ) in H that are mutually

140


orthogonal and have norm 1. A natural question arises whether every element of H can be represented (or at least approximated) by linear combinations of the (eα ). In L2 ([−π, π]) this leads, for example, to Fourier series representaP tions of functions in the form f (x) = ∞ n=0 (f, ψn )ψn , where the orthonormal functions are ψ0 = √12π , ψ2n (x) = √1π cos nx, ψ2n−1 (x) = √1π sin nx, and the series converges in L2 -norm. The completeness of L2 is crucial in ensuring the existence of such an orthonormal basis.

5.3 The Lp spaces: completeness More generally, the space Lp (E) is obtained when we integrate the pth powers of |f |. For p ≥ 1, we say that f ∈ Lp (and similarly for Lp (E) and Lp (E, C)) if |f |p is integrable (with the same convention of identifying f and g when they are a.e. equal). Some work will be required to check that Lp is a vector space and that the ‘natural’ generalization of the norm introduced for L2 is in fact a norm. We shall need p ≥ 1 to achieve this.

Definition 5.5 For each p ≥ 1, p < ∞, we define (identifying classes and functions) Z p L (E) = {f : |f |p dm is finite} E

p

and the norm on L is defined by kf kp =

Z

E

|f |p dm

p1

.

(With this in mind, we denoted the norm in L1 (E) by kf k1 and that in L2 (E) by kf k2.) Recall Definition 3.2: for any measurable function f : E → [0, ∞] ess supf := inf{c : |f | ≤ c a.e.}. More precisely, if F = {c ≥ 0 : m{|f |−1 ((c, ∞])} = 0}, we set ess supf = inf F (with the convention inf Ø = +∞). It is easy to see that the infimum belongs to F .


141

Definition 5.6 A measurable function f satisfying ess sup|f | < ∞ is said to be essentially bounded and the set of all essentially bounded functions on E is denoted by L∞ (E) (again with the usual identification of functions with a.e. equivalence classes), with the norm ||f ||∞ = ess supf. It is clear from Proposition 3.9 and the obvious identity ess sup(cf ) = c(ess supf ) that L∞ (E) is a vector space. We shall need to justify the notation by showing that for each p (1 ≤ p ≤ ∞), (Lp (E), k · k) is a normed vector space. First we observe that Lp (E) is a vector space for 1 ≤ p < ∞. If f and g belong to Lp , then they are measurable, hence so are cf and f + g. We have |cf (x)|p = |c|p |f (x)|p hence Z Z p1 p1 p p kcf kp = |cf (x)| dx = |c| |f (x)| dx = |c|kf kp . Next |f (x) + g(x)|p ≤ 2p max{|f (x)|p , |g(x)|p } and so kf + gkp is finite if only kf kp and kgkp are. Moreover, if kf kp = 0 then |f (x)|p = 0 almost everywhere and so f (x) = 0 almost everywhere. The converse is obvious. The triangle inequality kf + gkp ≤ kf kp + kgkp is by no means obvious for general p ≥ 1: we need to derive a famous inequality due to H¨ older, which is also extremely useful in many contexts, and generalizes the Schwarz inequality.

Remark 5.3 Before tackling this, we observe that the case of L∞ (E) is rather easier: |f +g| ≤ |f |+|g| at once implies that kf +gk∞ ≤ kf k∞ +kgk∞ and similarly |af | = |akf | gives kaf k∞ = kakkf k∞. Thus L∞ (E) is a vector space and since kf k∞ = 0 obviously holds if and only if f = 0 a.e., it follows that k · k∞ is a norm on L∞ (E).Exercises 3.6 and 5.6 show that this norm cannot be induced by any inner product on L∞ .

Lemma 5.7 For any non-negative real numbers x, y and all α, β ∈ (0, 1) with α + β = 1 we have xα y β ≤ αx + βy.

142


Proof If x = 0 the claim is obvious. So take x > 0. Consider f (t) = (1 − β) + βt − tβ for t ≥ 0 and β as given. We have f 0 (t) = β − βtβ−1 = β(1 − tβ−1 ) and since 0 < β < 1, f 0 (t) < 0 on (0,1). So f decreases on [0, 1] while f 0 (t) > 0 on (1, ∞), hence f increases on [1, ∞). So f (1) = 0 is the only minimum point of f on β [0, ∞), that is f (t) ≥ 0 for t ≥ 0. Now set t = xy , then (1 − β) + β xy − xy ≥ 0, β β that is, xy ≤ α + β xy . Writing x = xα+β we have xα+β xy ≤ αx + βx xy so that xα y β ≤ αx + βy as required.

Theorem 5.8 (Hölder’s Inequality) If

1 p

+

1 q

= 1, p > 1, then for f ∈ Lp (E), g ∈ Lq (E), we have f g ∈ L1 and kf gk1 ≤ kf kp kgkq

that is Z

|f g| dm ≤

Z

p

|f | dm

p1 Z

q

|g| dm

1q

.

Proof Step 1. Assume that kf kp = kgkq = 1, so we only need to show that kf gk1 ≤ 1. We apply Lemma 5.7 with α = 1p , β = 1q , x = |f |p , y = |g|q , then we have 1

1

|f g| = x p y q ≤

1 p 1 q |f | + |g| . p q

Integrating we obtain Z Z Z 1 1 1 1 p |f g| dm ≤ |f | dm + |g|q dm = + = 1 p q p q R R since |f |p dm = 1, |g|q dm = 1. So we have kf gk1 ≤ 1 as required.

Step 2. For general f ∈ Lp and g ∈ Lq we write kf kp = a, kgkq = b for some a, b > 0. (If either a or b is zero, then one of the functions is zero almost everywhere and the inequality is trivial.) Hence the functions f˜ = a1 f , g˜ = 1b g ˜ p k˜ satisfy the assumption of Step 1, and so kf˜g˜k1 ≤ kfk gkq . This yields 1 1 1 kf gk1 ≤ kf kp kgkq ab a b and after multiplying by ab the result is proved.


143

Letting p = q = 2 and recalling the definition of the scalar product in L2 we obtain the following now familiar special case of H¨ older’s inequality.

Corollary 5.9 (Schwarz Inequality) If f, g ∈ L2 , then

|(f, g)| ≤ kf k2 kgk2.

We may now complete the verification that k · kp is a norm on Lp (E).

Theorem 5.10 (Minkowski’s Inequality) For each p ≥ 1, f, g ∈ Lp (E) kf + gkp ≤ kf kp + kgkp .

Proof Assume 1 < p < ∞ (the case p = 1 was done earlier). We have |f + g|p = |(f + g)(f + g)p−1 | ≤ |f kf + g|p−1 + |gkf + g|p−1 , and also taking q such that

1 p

+

1 q

= 1, in other words, p + q = pq, we obtain

|f + g|(p−1)q = |f + g|p < ∞. Hence (f + g)p−1 ∈ Lq and k(f + g)

p−1

kq =

Z

p

|f + g| dm

q1

.

We may apply H¨ older’s inequality: Z Z Z p p−1 |f + g| dm ≤ |f kf + g| dm + |gkf + g|p−1 dm ≤

Z

+ =

|f |p dm

Z Z

p1 Z

|g|p dm p

|f | dm

|f + g|p dm

p1 Z

p1

+

q1

|f + g|p dm

Z

p

|g| dm

1q

p1 !

·A

144


R 1 with A = |f + g|p dm q . If A = 0 then kf + gkp = 0 and there is nothing to prove. So suppose A > 0 and divide by A kf + gkp =

Z

p

1− 1q

|f + g| dm Z 1 p |f + g| dm = A Z p1 Z p1 ≤ |f |p dm + |g|p dm = kf kp + kgkp

which was to be proved. Next we prove that Lp (E) is an example of a complete normed space (i.e. a Banach space) for 1 < p < ∞, i.e. that every Cauchy sequence in Lp (E) converges in norm to an element of Lp (E). We sometimes refer to convergence of sequences in the Lp −norm as convergence in pth mean. The proof is quite similar to the case p = 1.

Theorem 5.11 The space Lp (E) is complete for 1 < p < ∞.

Proof Given a Cauchy sequence fn (that is, kfn − fm kp → 0 as n, m → ∞) we find a subsequence fnk with 1 kfn − fnk kp < k 2 for all k ≥ 1 and we set gk =

k X i=1

|fni+1 − fni |,

g = lim gk = k→∞

∞ X i=1

|fni+1 − fni |.

Pk The triangle inequality yields kgk kp ≤ i=1 21i < 1 and we can apply Fatou’s lemma to the non-negative measurable functions gkp , k ≥ 1, so that Z Z p p kgkp = lim gk dm ≤ lim inf gkp dm ≤ 1. n→∞

k→∞

P Hence g is almost everywhere finite and fn1 + i≥1 (fni+1 − fni ) converges absolutely almost everywhere, defining a measurable function f as its sum.


145

We need to show that f ∈ Lp . Note first that f = limk→∞ fnk a.e., and given ε > 0 we can find N such that kfn − fm kp < ε for m, n ≥ N . Applying Fatou’s lemma to the sequence (|fni − fm |p )i≥1 , letting i → ∞, we have Z Z p |f − fm | dm ≤ lim inf |fni − fm |p dm ≤ εp . i→∞

Hence f − fm ∈ Lp and so f = fm + (f − fm ) ∈ Lp and we have kf − fm kp < ε for all m ≥ N . Thus fm → f in Lp -norm as required. The space L∞ (E) is also complete, since for any Cauchy sequence (fn ) in L (E) the union of the null sets where |fk (x)| > kf k∞ or |fn (x) − fm (x)| > kfn −fm k∞ for k, m, n ∈ N, is still a null set, F say. Outside F the sequence (fn ) converges uniformly to a bounded function, f say. It is clear that kfn −f k∞ → 0 and f ∈ L∞ (E), so we are done. ∞

Exercise 5.8 Is the sequence 1 gn (x) = 1(0, n1 ] (x) √ x Cauchy in L4 ? We have the following relations between the Lp spaces for different p which generalize Proposition 5.3.

Theorem 5.12 If E has finite Lebesgue measure, then Lq (E) ⊆ Lp (E) when 1 ≤ p ≤ q ≤ ∞.

Proof Note that |f (x)|p ≤ 1 if |f (x)| ≤ 1. If |f (x)| ≥ 1, then |f (x)|p ≤ |f (x)|q . Hence |f (x)|p ≤ 1 + |f (x)|q , Z Z Z Z |f |p dm ≤ 1 dm + |f |q dm = m(E) + |f |q dm < ∞, E

so if m(E) and

E

R

q

E

E

E

|f | dm are finite, the same is true for

R

E

|f |p dm.

146


5.4 Probability 5.4.1 Moments Random variables belonging to spaces Lp (Ω), where the exponent p ∈ N, play an important role in probability.

Definition 5.7 The moment of order n of a random variable X ∈ Ln (Ω) is the number E(X n ),

n = 1, 2, . . .

Write E(X) = µ; then central moments are given by E(X − µ)n ,

n = 1, 2, . . .

By Theorem 4.28 moments are determined by the probability distribution: Z n E(X ) = xn dPX (x), n

E((X − µ) ) =

Z

(x − µ)n dPX (x),

Z

(x − µ)n fX (x) dx.

and if X has a density fX then by Theorem 4.32 we have Z E(X n ) = xn fX (x) dx, E((X − µ)n ) =

Proposition 5.13 If E(X n ) is finite for some n, then for k ≤ n, E(X k ) are finite. If E(X n ) is infinite, then the same is true for E(X k ) for k ≥ n. Hint Use Theorem 5.12.

Exercise 5.9 Find X so that E(X 2 ) = ∞, E(X) < ∞. Can such an X have E(X) = 0? Hint You may use some previous examples in this chapter.


147

Definition 5.8 The variance of a random variable is the central moment of second order: Var(X) = E(X − E(X))2 . Clearly, writing µ = E(X), Var(X) = E(X 2 − 2µX + µ2 ) = E(X 2 ) − 2µE(X) + µ2 = E(X 2 ) − µ2 . This shows that the first two moments determine the second central moment. This may be generalized to arbitrary order and what is more, this relationship also goes the other way round.

Proposition 5.14 Central moments of order n are determined by moments of order k for k ≤ n. Hint Use the binomial theorem and linearity of the integral.

Proposition 5.15 Moments of order n are determined by central moments of order k for k ≤ n. Hint Write E(X n ) as E((X − µ + µ)n ) and then use the binomial theorem.

Exercise 5.10 Find Var(aX) in terms of Var(X).

Example 5.5 If X has the uniform distribution on [a, b], that is, fX (x) = Z

xfX (x) dx =

1 b−a

Z

b

x dx = a

1 b−a 1[a,b] (x)

1 1 2b 1 x | = (a + b). b−a2 a 2

Exercise 5.11 Show that for uniformly distributed X, VarX =

1 12 (b

− a)2 .

then

148


Exercise 5.12 Find the variance of (a) a constant random variable X, X(ω) = a for all ω (b) X : [0, 1] → R given by X(ω) = min{ω, 1 − ω} (the distance to the nearest endpoint of the interval [0, 1]) (c) X : [0, 1]2 → R, the distance to the nearest edge of the square [0, 1]2 . We shall see that for the Gaussian distribution the first two moments determine the remaining ones. First we compute the expectation:

Theorem 5.16 √

1 2πσ

Z

xe−

(x−µ)2 2σ2

dx = µ.

R

Proof R R Make the substitution z = x−µ for R , σ , then, writing Z Z Z (x−µ)2 1 σ z2 µ z2 √ ze− 2 dz + √ e− 2 dz. xe− 2σ2 dx = √ 2πσ 2π 2π Notice that the first √ integral is zero since the integrand is an odd function. The second integral is 2π, hence the result. So the parameter µ in the density is the mathematical expectation. We show now that σ 2 is the variance.

Theorem 5.17 √

1 2πσ

Z

R

(x − µ)2 e−

(x−µ)2 2σ2

dx = σ 2 .

Proof Make the same substitution as before: z = x−µ σ ; then Z Z (x−µ)2 1 σ2 z2 2 − 2σ2 √ (x − µ) e dx = √ z 2 e− 2 dz. 2πσ 2π


149

2

Integrate by parts u = z, v = ze−z /2 , to get Z Z σ2 σ2 σ2 z 2 +∞ z2 z2 √ +√ z 2 e− 2 dz = − √ ze− 2 e− 2 dz = σ 2 −∞ 2π 2π 2π since the first term vanishes.

Note that the odd central moments for a Gaussian random variable are zero: the integrals Z (x−µ)2 1 √ (x − µ)2k+1 e− 2σ2 dx 2πσ vanish since after the above substitution we integrate an odd function. By repeating the integration by parts argument one can prove that E(X − µ)2k = 1 · 3 · 5 · · · (2k − 1)σ k .

Example 5.6 1 Let us consider the Cauchy density π1 1+x 2 and try to compute the expectation (we shall see it is impossible): Z 1 +∞ x 1 2 2 dx = lim ln(1 + xn ) − lim ln(1 + yn ) yn →−∞ π −∞ 1 + x2 2π xn →+∞

for some sequences xn , yn . The result, if finite, should not depend on their choice, however if we set for example xn = ayn , then we have 1 + ayn2 2 2 lim ln(1 + xn ) − lim ln(1 + yn ) = lim ln = ln a xn →+∞ yn →−∞ yn →∞ 1 + yn2 which is a contradiction. As a consequence, we see that for the Cauchy density the moments do not exist.

Remark 5.4 We give without proof a simple relation between the characteristic function and the moments: (Recall that ϕX (t) = E(eitX ) – see Definition 4.5.) If ϕX is k-times continuously differentiable then X has finite kth moment and 1 dk E(X k ) = k k ϕX (0). i dt Conversely, if X has kth moment finite then ϕX (t) is k-times differentiable and the above formula holds.

150


5.4.2 Independence The expectation provides a useful criterion for the independence of two random variables.

Theorem 5.18 The random variables X, Y are independent if and only if E(f (X)g(Y )) = E(f (X))E(g(Y ))

(5.4)

holds for all Borel measurable bounded functions f , g.

Proof Suppose that (5.4) holds and take any Borel sets B1 , B2 . Let f = 1B1 , g = 1B2 and application of (5.4) gives Z Z Z 1B1 (X(ω))1B2 (Y (ω)) dP (ω) = 1B1 (X(ω)) dP (ω) 1B2 (Y (ω)) dP (ω). Ω

Ω

Ω

The left-hand side equals Z 1B1 ×B2 (X(ω), Y (ω)) dP (ω) = P ((X ∈ B1 ) ∩ (Y ∈ B2 )), Ω

whereas the right-hand side is P (X ∈ B1 )P (Y ∈ B2 ), thus proving the independence of X and Y . Suppose now that X, Y are independent. Then (5.4) holds for f = 1B1 , g = 1B2 , B1 , B2 Borel sets, by the above argument. By linearity we extend the P P formula to simple functions: ϕ = b i 1B i , ψ = j c j 1C j , X X E(ϕ(X)ψ(Y )) = E( bi 1Bi (X) cj 1Cj (Y )) X = bi cj E(1Bi (X)1Cj (Y )) i,j

=

X

bi cj E(1Bi (X))E(1Cj (Y ))

i,j

=

X i

bi E(1Bi (X))

X

cj E(1Cj (Y ))

j

= E(ϕ(X))E(ψ(Y )).

We approximate general f , g by simple functions and the dominated convergence theorem (f , g are bounded) extends the formula to f , g.


151

Proposition 5.19 Assume that X, Y are independent random variables. Show that if E(X) = 0, E(Y ) = 0, then E(XY ) = 0. Hint The above theorem cannot be applied with f (x) = x, g(x) = x (these functions are not bounded). So some approximation is required. The expectation is nothing but an integral so the number (X, Y ) = E(XY ) is the inner product in the space L2 (Ω) of random variables square integrable with respect to P . Hence independence implies orthogonality in this space. If the expectation of a random variable is non-zero, we modify the notion of orthogonality. The idea is that adding (or subtracting) a number does not destroy or improve independence.

Definition 5.9 For a random variable with finite µ = E(X) we write Xc = X − E(X) and we call Xc a centred random variable (clearly E(Xc ) = 0). The covariance of X and Y is defined as Cov(X, Y ) = (Xc , Yc ) = E (X − E(X))(Y − E(Y )) . The correlation is the cosine of the angle between Xc and Yc : ρX,Y =

(Xc , Yc ) Cov(X, Y ) = . kXk2kY k2 kXk2kY k2

We say that X, Y are uncorrelated if ρ = 0. Note that some elementary algebra gives a more convenient expression for the covariance: Cov(X, Y ) = E(XY ) − E(X)E(Y ). Thus uncorrelated X, Y satisfy E(XY ) = E(X)E(Y ). Clearly independent random variables are uncorrelated; it is sufficient to take f (x) = x − E(X), g(x) = x − E(Y ) in Theorem 5.18. The converse is not true in general, although – as we shall see in Chapter 6 – it holds for Gaussian random variables.

Example 5.7 Let Ω = [−1, 1] with Lebesgue measure: P = 12 m|[−1,1] , X = x, Y = x2 . R1 Then E(X) = 0, E(XY ) = −1 x3 dx = 0, hence Cov(X, Y ) = 0 and thus ρX,Y = 0. However X, Y are not independent. Intuitively this is clear since

152


Y = X 2 , so that each of X, Y is a function of the other. Specifically, take A = B = [− 21 , 12 ] and compare (as required by Definition 3.3) the probabilities P (X −1 (A) ∩ Y −1 (A)) and P (X −1 (A))P (Y −1 (A)). We obtain X −1 (A) = A, Y −1 (A) = [− √12 , √12 ], hence P (X −1 (A) ∩ Y −1 (A)) = √12 , P (X −1 (A)) = 12 , P (Y −1 (A)) = √12 and so X and Y are not independent.

Exercise 5.13 Find the correlation ρX,Y if X = 2Y + 1.

Exercise 5.14 Take Ω = [0, 1] with Lebesgue measure and let X(ω) = sin 2πω, Y (ω) = cos 2πω. Show that X, Y are uncorrelated but not independent. We close the section with two further applications.

Proposition 5.20 The variance of the sum of uncorrelated random variables is the sum of their variances: n n X X Var( Xi ) = Var(Xi ). i=1

i=1

Hint To avoid cumbersome notation first prove the formula for two random variables Var(X + Y ) = Var(X) + Var(Y ) using the formula Var(X) = E(X 2 ) − (EX)2 .

Proposition 5.21 Suppose that X, Y are independent random variables. Then we have the following formula for the characteristic function: ϕX+Y (t) = ϕX (t)ϕY (t). More generally, if X1 , . . . , Xn are independent, then ϕX1 +···+Xn (t) = ϕX1 (t) · · · ϕXn (t). Hint Use the definition of characteristic functions and Theorem 5.18.


153

5.4.3 Conditional Expectation (first construction) The construction of orthogonal projections in complete inner product spaces, undertaken in Section 5.2.3, allows us to provide a preview of perhaps the most important concept in modern probability theory: the conditional expectation of an F-measurable integrable random variable X, given a σ-field G contained in F (i.e. such that every set in G also belongs to F). We study this idea in detail in Chapter 7 where we will also justify the definition below by reference to more familiar concepts, but the construction of the conditional expectation as a G-measurable random variable can be achieved for any integrable X with the tools we have readily to hand. Our argument owes much to the elegant construction given in [12].

Definition 5.10 Let (Ω, F, P ) be a probability space and suppose that G is a sub-σ-field of F. 1 1 1 GivenR X ∈ L1 (Ω, R F, P ) = L (F) there exists Y ∈ L (Ω, G, P ) = L (G) such that G Y dP = G X dP for every G ∈ G. We write Y = E(X|G) and call Y the conditional expectation of X given G. These conditions define Y uniquely up to P -null sets. Theorem 4.15, applied to P instead of m and GR instead of RM, implies that Y is P -a.s. unique: if Z ∈ L1 (G) also satisfies G Z dP = G X dP for every G ∈ G, then Z = Y P -a.s. This is often expressed by saying that Y is a version of E(X|G ): by definition of L1 (Ω, G, P ) the uniqueness claim is that all versions belong to the same equivalence class in L1 (Ω, G, P ) under the equivalence relation f ≡ g if and only if P ({ω ∈ Ω : f (ω) 6= g(ω)}). In accordance with our convention (Section 5.1) we shall nonetheless continue to work with functions rather than with equivalence classes. To construct Y we first restrict attention to the case when X ∈ L2 (F) = L2 (Ω, F, P ). By Theorem 5.11, the inner product space H = L2 (F) is complete, and the vector subspace K = L2 (G) is a complete subspace. Thus the construction of the orthogonal projection in Section 5.2.3 applies, and provides an element of L2 (G), which by our convention we represent as a function Y ∈ L2 (G), such that (X − Y ) is orthogonal to K. By definition of the inner product in L2 this means that Z (X − Y, Z) = (X − Y )Z dP = 0 Ω

2 2 Rfor every Z R∈ L (G). In particular, since 1G ∈ L (G) for every G ∈ G, we have Y dP = G X dP. G

154


To construct Y for an arbitrary X ∈ L1 (F)) we proceed in four stages: (i) first, note that if the result has been proved for non-negative functions + − in L1 (F)), we can consider X = X +R− X − . By hypothesis there are R R Y −, Y ∈ 1 + + L such that for G ∈ G both G Y dP = RG X dP and R (G)) R G Y dP = − X dP. Subtracting on both sides we obtain G Y dP = G X dP , where G Y = Y + − Y − ∈ L1 (G). So we need only verify the result for non-negative integrable X. (ii) second, if a random variable Z is bounded, it is in L2 (F) by Theorem 5.12, since P (Ω) is finite. R Hence Z has R a conditional expectation, i.e. there exists W ∈ L2 (G) such that G W dP = G Z dP for G ∈ G. Also: if Z ≥ 0, then W ≥ 0 P -a.s. To see this, suppose that W takes negative values with positive probability. ThenR there exists n ≥ 1 such that the set G = {W < − n1 } ∈ G has R R P (G) > 0. Thus G W dP < − n1 P (G) < 0. But G W dP = G Z dP ≥ 0, since Z ≥ 0 (Proposition 4.11). The contradiction shows that W ≥ 0, P -a.s. (iii) Now take an arbitrary X ≥ 0 in L1 (F), and for n ≥ 1 set Xn = min(X, n). Then Xn is bounded and non-negative, so part (ii) applies to Xn , R R yielding a non-negative Yn ∈ L2 (G) with G Yn dP = G Xn dP. Since (Xn ) is increasing with n, for any fixed n the bounded random variable Z = Xn+1 −Xn is non-negative, and has a conditional expectation W ≥ 0, as in (ii). But so have both Xn+1 and Xn ,and the a.s. uniqueness property therefore implies that W = Yn+1 − Yn P -a.s. Therefore (Yn ) also increases (a.s.) with n. (iv) Finally, set Y (ω) = lim supn→∞ Yn (ω) for each ω ∈ Ω. By Theorem 3.5 Y is G-measurable, and the Rsequence (YnR) increases a.s. R to Y. Moreover 0 ≤ X ∈ L1 (F) and for G ∈ G, G Yn dP = G Xn dP ≤ G X dP n and as n → ∞, fn % f . The sections of simple measurable functions are simple and measurable. This is clear for the indicator functions as observed above, and next we use the fact that the section of the sum is the sum of the sections. Finally, it is clear that the sections of fn converge to the sections of f and since measurability is preserved in the limit, the theorem is proved.

6. Product measures

171

Corollary 6.8 The functions ω1 7→

Z

f (ω1 , ω2 ) dP2 (ω2 ), Ω2

ω2 7→

Z

f (ω1 , ω2 ) dP1 (ω1 ) Ω1

are F1 , F2 -measurable, respectively.

Proof The integrals may be taken (being possibly infinite) due to measurability of the functions in question. By the monotone convergenceRtheorem, they are limRits of the integrals of the sections of fn . The integrals Ω1 fn (ω1 , ω2 ) dP1 (ω1 ), Ω2 fn (ω1 , ω2 ) dP2 (ω2 ) are simple functions, and hence the limits are measurable.

Theorem 6.9 Let f be a measurable non-negative function defined on Ω1 × Ω2 . Then Z Z Z f (ω1 , ω2 ) d(P1 × P2 )(ω1 , ω2 ) = f (ω1 , ω2 ) dP2 (ω2 ) dP1 (ω1 ) Ω1 ×Ω2

Ω1

=

Z

Ω2

Z

Ω2

f (ω1 , ω2 ) dP1 (ω1 ) Ω1

dP2 (ω2 ).

(6.4)

Proof For the indicator function of a rectangle A1 ×A2 each side of (6.4) just becomes P1 (A1 )P2 (A2 ). Then by additivity of the integral the formula is true for simple functions. Monotone approximation of any measurable f by simple functions allows us to extend this formula to the general case.

Theorem 6.10 (Fubini’s Theorem) If f ∈ L1 (Ω1 × Ω2 ) then the sections are integrable in appropriate spaces, the functions Z Z ω1 7→ f (ω1 , ω2 ) dP2 (ω2 ), ω2 7→ f (ω1 , ω2 ) dP1 (ω1 ) Ω2

1

Ω1

1

are in L (Ω1 ), L (Ω2 ), respectively, and (6.4) holds: in concise form it reads Z Z Z Z Z f d(P1 × P2 ) = f dP2 dP1 = f dP1 dP2 . Ω1 ×Ω2

Ω1

Ω2

Ω2

Ω1

172


Proof This relation is immediate by the decomposition f = f + − f − and the result proved for non-negative functions. The integrals are finite since if f ∈ L1 then f + , f − ∈ L1 and all the integrals on the right are finite.

Remark 6.2 The whole procedure may be extended to the product of an arbitrary finite number of spaces. In particular, we have a method of constructing ndimensional Lebesgue measure as the completion of the product of n copies of one-dimensional Lebesgue measure.

Example 6.1 Let Ω1 = Ω2 = [0, 1], P1 = P2 = m[0,1] , 1 if 0 < y < x < 1 x2 f (x, y) = 0 otherwise. We shall see that the integral of f over the square is infinite. For this we take a non-negative simple function dominated by f and compute its integral. Let 1 ϕ(x, y) = n if f (x, y) ∈ [n, n + 1). Then ϕ(x, y) = n if x > y, x ∈ ( √n+1 , √1n ]. 1 1 1 The area of this set is 2 ( n − n+1 ) and Z ∞ ∞ X X 1 1 1 1 1 ϕ dm2 = n ( − )= = ∞. 2 n n + 1 2 n+1 [0,1]2 n=1 n=1 Hence the function

g(x, y) =

 

1 x2

− 12  y 0

if 0 < y < x < 1 if 0 < x < y < 1 otherwise

is not integrable since the integral of g + is infinite (the same is true for the integral of g − ).

Exercise 6.1 For g from the above example show that Z 1Z 1 Z 1Z g(x, y) dx dy = −1, 0

0

0

1

g(x, y) dy dx = 1 0

which shows that the iterated integrals may not be equal if Fubini’s theorem condition is violated.

6. Product measures

173

The following proposition opens the way for many applications of product measures and Fubini’s theorem.

Proposition 6.11 Let f : R → R be measurable and positive. Consider the set of all points in the upper half-plane being below the graph of f : Af = {(x, y) : 0 ≤ y < f (x)}. R Show that Af is m2 -measurable and m2 (Af ) = f (x) dx.

Hint For measurability ‘fill’ Af with rectangles using the approximation of f by simple functions. Then apply the definition of the product measure.

Exercise 6.2 Compute

R

[0,3]×[−1,2]

x2 y dm2 .

Exercise 6.3 Compute the area of the region inside the ellipse

x2 a2

+

y2 b2

= 1.

6.5 Probability 6.5.1 Joint distributions Let X, Y be two random variables defined on the same probability space (Ω, F, P ). Consider the random vector (X, Y ) : Ω → R2 . Its distribution is the measure defined for the Borel sets on the plane given by P(X,Y ) (B) = P ((X, Y ) ∈ B),

B ⊂ R2 .

If this measure can be written as P(X,Y ) (B) =

Z

f(X,Y ) (x, y) dm2 (x, y) B

for some integrable f(X,Y ) , then we say that X, Y have a joint density.

174


The joint distribution determines the distributions of one-dimensional random variables X, Y : PX (A) = P(X,Y ) (A × R), PY (A) = P(X,Y ) (R × A), for Borel A ⊂ R, these are called marginal distributions. If X, Y have a joint density, then both X and Y are absolutely continuous with densities given by Z fX (x) = f(X,Y ) (x, y) dy, R

fY (y) =

Z

f(X,Y ) (x, y) dx. R

The following example shows that the converse is not true in general.

Example 6.2 Let Ω = [0, 1] with P = m[0,1] and let (X, Y )(ω) = (ω, ω). This vector does not have density since P(X,Y ) ({(x, y) : x = y}) = 1 and for any integrable function R 2 f : R → R, {(x,y):x=y} f (x, y) dm2 (x, y) = 0; a contradiction. However the marginal distributions PX , PY are absolutely continuous with the densities fX = fY = 1[0,1] .

Example 6.3 1 1A , with Borel A simple example of joint density is the uniform one: f = m(A) 2 A ⊂ R . A particular case is A = [0, 1] × [0, 1], then clearly the marginal densities are 1[0,1] .

Exercise 6.4 Take A to be the square with corners at (0, 1), (1, 0), (2, 1), (1, 2). Find the marginal densities of f = 1A .

Exercise 6.5 1 (x2 + y 2 ) if 0 < x < 2, 1 < y < 4 and zero otherwise. Let fX,Y (x, y) = 50 Find P (X + Y > 4), P (Y > X).

The two-dimensional Gaussian (normal) density is given by n(x, y) =

1 2π

p

1 − ρ2

exp −

1 (x2 − 2ρxy + y 2 ) . 2 2(1 − ρ )

(6.5)

6. Product measures

175

It can be shown that ρ is the correlation of X, Y , random variables whose densities are the marginal densities of n(x, y), (see [9]). Joint densities enable us to compute the distributions of various functions of random variables. Here is an important example.

Theorem 6.12 If X, Y have joint density fX,Y , then the density of their sum is given by Z fX+Y (z) = fX,Y (x, z − x) dx. (6.6) R

Proof We employ the distribution function: FX+Y (z) = P (X + Y ≤ z)

= PX,Y ({(x, y) : x + y ≤ z}) Z Z = fX,Y (x, y) dxdy {(x,y):x+y≤z} z−x

= =

Z Z

fX,Y (x, y) dydx

R −∞ Z z Z −∞

R

fX,Y (x, y 0 − x) dxdy 0

(we have used the substitution y 0 = y + x and Fubini’s theorem), which by differentiation gives the result.

Exercise 6.6 Find fX+Y if fX,Y = 1[0,1]×[0,1] .

6.5.2 Independence again Suppose that the random variables X, Y are independent. Then for a Borel rectangle: B = B1 × B2 we have P(X,Y ) (B1 × B2 ) = P ((X, Y ) ∈ B1 × B2 )

= P ((X ∈ B1 ) ∩ (Y ∈ B2 )) = P (X ∈ B1 )P (Y ∈ B2 ) = PX (B1 )PY (B2 )

176


and so the distribution P(X,Y ) coincides with the product measure PX × PY on rectangles, therefore they are the same. The converse is also true:

Theorem 6.13 The random variables X, Y are independent if and only if P(X,Y ) = PX × PY .

Proof The ‘only if’ part was shown above. Suppose that P(X,Y ) = PX × PY and take any Borel sets B1 , B2 . The same computation shows that P ((X ∈ B1 ) ∩ (Y ∈ B2 )) = P (X ∈ B1 )P (Y ∈ B2 ), i.e. X and Y are independent. We have a useful version of this theorem in the case of absolutely continuous random variables.

Theorem 6.14 If X, Y have a joint density, then they are independent if and only if f(X,Y ) (x, y) = fX (x)fY (y).

(6.7)

If X and Y are absolutely continuous and independent, then they have a joint density and it is given by (6.7).

Proof Suppose f(X,Y ) is the joint density of X, Y . If they are independent, then Z f(X,Y ) (x, y) dm2 (x, y) = P ((X, Y ) ∈ B1 × B2 ) B1 ×B2

= P (X ∈ B1 )P (Y ∈ B2 ) Z Z = fX (x) dm(x) fY (y) dm(y) B1 B2 Z = fX (x)fY (y) dm2 (x, y) B1 ×B2

6. Product measures

177

which implies (6.7). The same computation shows the converse: Z P ((X, Y ) ∈ B1 × B2 ) = f(X,Y ) (x, y) dm2 (x, y) ZB1 ×B2 = fX (x)fY (y) dm2 (x, y) B1 ×B2

= P (X ∈ B1 )P (Y ∈ B2 ).

For the final claim note that the function fX (x)fY (y) plays the role of the joint density if X and Y are independent.

Corollary 6.15 If Gaussian random variables are orthogonal, then they are independent.

Proof Inserting ρ = 0 into (6.5) we immediately see that the two-dimensional Gaussian density is the product of the one-dimensional ones.

Proposition 6.16 The density of the sum of independent random variables with densities fX , fY is given by Z fX+Y (z) = fX (x)fY (z − x) dx. R

Exercise 6.7 Suppose that the joint density of X, Y is 1A where A is the square with corners at (0, 1), (1, 0), (2, 1), (1, 2). Are X, Y independent?

Exercise 6.8 Find P (Y > X) and P (X + Y > 1), if X, Y are independent with fX = 1[0,1] , fY = 21 1[0,2] .

178


6.5.3 Conditional probability We consider the case of two random variables X, Y with joint density fX,Y (x, y). Given Borel sets A, B, we compute P (X ∈ A, Y ∈ B) P (X ∈ A) R f(X,Y ) (x, y) dm2 (x, y) = A×B R f (x) dm(x) A X Z R f(X,Y ) (x, y) dx AR = dy f (x) dx B A X

P (Y ∈ B|X ∈ A) =

using Fubini’s theorem. So the conditional distribution of Y given X ∈ A has a density R f(X,Y ) (x, y) dx h(y|X ∈ A) = A R . A fX (x) dx

The case where A = {a} does not make sense here since then we would have zero in the denominator. However, formally we may put h(y|X = a) =

f(X,Y ) (a, y) fX (a)

which makes sense if only fX (a) 6= 0. This restriction turns out to be not relevant since Z P ((X, Y ) ∈ {(x, y) : fX (x) = 0}) = f(X,Y ) (x, y) dx dy {(x,y):fX (x)=0} Z Z = f(X,Y ) (x, y) dy dx {x:fX (x)=0} R Z = fX (x) dx {x:fX (x)=0}

= 0.

We may thus define the conditional probability of Y ∈ B given X = a by means of h(y|X = a) which we briefly write as h(y|a): Z P (Y ∈ B|X = a) = h(y|a) dy B

and the conditional expectation E(Y |X = a) =

Z

yh(y|a) dy. R

6. Product measures

179

This can be viewed as a random variable with X as the source of randomness. Namely, for ω ∈ Ω we write Z E(Y |X)(ω) = yh(y|X(ω)) dy. R

This function is of course measurable with respect to the σ-field generated by X. The expectation of this random variable can be computed using Fubini’s theorem: Z E(E(Y |X)) = E( yh(y|X(ω)) dy) Z ZR = yh(y|x) dyfX (x) dx R R Z Z = yf(X,Y ) (x, y) dx dy ZR R = yfY (y) dy R

= E(Y ).

More generally, for A ⊂ Ω, A = X −1 (B), B Borel, Z Z E(Y |X) dP = 1B (X)E(Y |X) dP A ZΩ Z = 1B (X(ω))( yh(y|X(ω)) dy) dP (ω) R ZΩ Z = 1B (x)yh(y|x) dyfX (x) dx ZR ZR = 1B (x)yf(X,Y ) (x, y) dx dy ZR R = 1A (X)Y dP ZΩ = Y dP. A

This provides a motivation for a general notion of conditional expectation of a random variable Y given random variable X: E(Y |X) is a random variable measurable with respect to the σ-field FX generated by X and such that for all A ∈ FX Z Z A

E(Y |X) dP =

Y dP.

A

We will pursue these ideas further in the next chapter.

180


Exercise 6.9 Let fX,Y = 1A , where A is the triangle with corners at (0, 0), (2, 0), (0, 1). Find the conditional density h(y, x) and conditional expectation E(Y |X = 1).

Exercise 6.10 Let fX,Y (x, y) = (x + y)1A , where A = [0, 1] × [0, 1]. Find E(X|Y = y) for each y ∈ R.

6.5.4 Characteristic functions determine distributions We have now sufficient tools to prove a fundamental property of characteristic functions.

Theorem 6.17 (Inversion Formula) If the cumulative distribution function of a random variable X is continuous at a, b ∈ R, then Z c −iua 1 e − e−iub FX (b) − FX (a) = lim ϕX (u) du. c→∞ 2π −c iu

Proof First, by the definition of ϕX , 1 2π

Z

c −c

e−iua − e−iub 1 ϕX (u) du = iu 2π

Z

c

e−iua − e−iub iu

−c

Z

eiux dPX (x) du. R

We may apply Fubini’s theorem since |

e−iua − e−iub iux e |=| iu

Z

b

eiux d(x)| ≤ b − a

a

which is integrable with respect to PX × m[−c,c]. We compute the integral in u 1 2π 1 2π

Z

c −c

Z

c −c

e−iua − e−iub iux 1 e du = iu 2π

Z

sin u(x − a) − sin u(x − b) 1 du + u 2π

c −c

Z

eiu(x−a) − eiu(x−b) du = iu

c −c

cos u(x − a) − cos u(x − b) du. iu

6. Product measures

181

The second integral vanishes since the integrand is an odd function. We change variables in the first: y = u(x − a), z = u(x − b) and then it takes the form Z c(x−a) Z c(x−b) 1 sin y 1 sin z I(x, c) = dy − dz = I1 (x, c) − I2 (x, c), 2π −c(x−a) y 2π −c(x−b) z say. We employ the following elementary fact without proof: Z t sin y dy → π as t → ∞, s → −∞. y s Consider the following cases: 1. x < a, then also x < b and c(x−a) → −∞, c(x−b) → −∞, −c(x−a) → ∞, −c(x − b) → ∞ as c → ∞. Hence I1 (x, c) → − 21 , I2 (x, c) → − 12 and so I(x, c) → 0.

2. x > b, then also x > a, and c(x − a) → ∞, c(x − b) → ∞, −c(x − a) → −∞, −c(x − b) → −∞, as c → ∞ so I1 (x, c) → 12 , I2 (x, c) → 21 and the result is the same as in 1. 3. a < x < b hence I1 (x, c) → expression is 1.

1 2,

I2 (x, c) → − 21 and the limit of the whole

Write f (x) = limc→∞ I(x, c) (we have not discussed the values x = a, x = b but they are irrelevant as will be seen). Z c −iua Z 1 e − e−iub lim ϕX (u) du = lim I(x, c) dPX (x) c→∞ 2π −c c→∞ R iu Z = f (x) dPX (x) R

by Lebesgue’s dominated convergence theorem. The integral of f can be easily computed since f is a simple function: Z f (x) dPX (x) = PX ((a, b]) = FX (b) − FX (a) R

(PX ({a}) = PX ({b}) = 0 since FX is continuous at a and b).

Corollary 6.18 The characteristic function determines the probability distribution.

Proof Since FX is monotone, it is continuous except (possibly) at countably many points where it is right-continuous. Its values at discontinuity points can be

182


approximated from above by the values at continuity points. The latter are determined by the characteristic function via the inversion formula. Finally, we see that FX determines the measure PX . This is certainly so for B = (a, b]: PX ((a, b]) = FX (b) − FX (a). Next we show the same for any interval, then for finite unions of intervals, and the final extension to any Borel set is via the monotone class theorem.

Theorem 6.19 If ϕX is integrable, then X has a density which is given by Z ∞ 1 fX (x) = e−iux ϕX (u) du. 2π −∞

Proof The function f is well-defined. To show that it is a density of X we first show that it gives the right values of the probability distribution of intervals (a, b] where FX is continuous: ! Z b Z ∞ Z b 1 −iux fX (x) dx = ϕX (u) e dx du 2π −∞ a a ! Z c Z b 1 −iux = lim ϕX (u) e dx du c→∞ 2π −c a Z c 1 e−iua − e−iub = lim ϕX (u) du c→∞ 2π −c iu = FX (b) − FX (a) by the inversion formula. This extends to all a, b since FX is right continuous and the integral on the left R b is continuous with respect to a and b. Moreover, FX is non-decreasing so a fX (x) dx ≥ 0 for all a ≤ b hence fX ≥ 0. Finally Z ∞ fX (x) dx = lim FX (b) − lim FX (a) = 1 −∞

b→∞

a→−∞

so fX is a density.

6.5.5 Application to mathematical finance Classical portfolio theory is concerned with an analysis of the balance between risk and return. This balance is of fundamental importance, particularly in

6. Product measures

183

corporate finance, where the key concept is the cost of capital, which is a rate of return based on the level of risk of an investment. In probabilistic terms, return is represented by the expectation and risk by the variance. A theory which deals only with two moments of a random variable is relevant if we assume the normal (Gaussian) distribution of random variables in question, since in that case these two moments determine the distribution uniquely. We give a brief account of basic facts of portfolio theory under this assumption. Let k be a return on some investment in single period, that is, k(ω) = V (1,ω)−V (0) where V (0) is the known amount invested at the beginning, and V (0) V (1) is the random terminal value. A typical example which should be kept in mind is buying and selling one share of some stock. With a number of stocks i (0) available, we are facing a sequence ki of returns on stock Si , ki = Si (1,ω)−S , Si (0) but for simplicity we restrict our attention to just two, k1 , k2 . A portfolio is formed by deciding the percentage split, between holdings in S1 and S2 , of the initial wealth V (0) by choosing the weights w = (w1 , w2 ), w1 + w2 = 1. Then, (0) as is well known and elementary to verify, the portfolio of n1 = wS11V(0) shares of stock number one and n2 =

w2 V (0) S2 (0)

shares of stock number two, has return

kw = w 1 k1 + w 2 k2 . We assume that the vector (k1 , k2 ) is jointly normal with correlation coefficient ρ. We denote the expectations and variances of the ingredients by µi = E(ki ), σi2 = Var(ki ). It is convenient to introduce the following matrix c11 c12 σ12 ρσ1 σ2 C= = c21 c22 ρσ1 σ2 σ22 where c12 = c21 is the covariance between k1 and k2 . Assume (which is not elegant to do but saves us an algebraic detour) that C is invertible, with C −1 = [dij ]. By definition, the joint density has the form f (x1 , x2 ) =

2 1 1 X √ dij (xi − µi )(xj − µj )} exp{− 2 i,j=1 2π det C

It is easy to see that (6.5) is a particular case of this formula with µi = 0, σi = 1, −1 < ρ < 1. It is well known that the characteristic function ϕ(t1 , t2 ) = E(exp{i(t1 k1 + t2 k2 )}) of the vector (k1 , k2 ) is of the form ϕ(t1 , t2 ) = exp{i

2 X i=1

t i µi −

2 1 X cij ti tj }. 2 i.j=1

(6.8)

We shall show that the return on the portfolio is also normally distributed and we shall find the expectation and standard deviation. This can all be done in one step

184


Theorem 6.20 The characteristic ϕw function of kw is of the form 1 ϕw (t) = exp{it(w1 µ1 + w2 µ2 )) − t2 (w12 σ12 + w22 σ22 + 2w1 w2 ρσ1 σ2 )} 2

Proof By definition ϕw (t) = E(exp{itkw }) and using the form of kw we have ϕw (t) = E(exp{it(w1 k1 + w2 k2 )} = E(exp{itw1 k1 + itw2 k2 }) = ϕ(tw1 , tw2 )

by the definition of the characteristic function of a vector. Since the vector is normal, (6.8) immediately gives the result. The multi-dimensional version of Corollary 6.18 (which is easy to believe after mastering the one-dimensional case, but slightly tedious to prove, so we take it for granted, referring the reader to any probability textbook) shows that kw has normal distribution with µw = w 1 µ1 + w 2 µ2 2 σw = w12 σ12 + w22 σ22 + 2w1 w2 ρσ1 σ2

The fact that the variance of a portfolio can be lower than the variances of the components is crucial. These formulae are valid in general case (i.e. without the assumption of a normal distribution) and can be easily proved using the formula for kw . The main goal of this section was to see that the portfolio return is normally distributed.

Example 6.4 Suppose that the second component is not random, i.e. S2 (1) is a constant independent of ω. Then the return k2 is risk-free and it is denoted by r (the notation is usually reserved for the case where the length of the period is one year). It can be thought of as a bank account and it is convenient to assume that S2 (0) = 1. Then the portfolio of n shares purchased at the price S1 (0) and m units of the bank account has the value V (1) = nS1 (1) + m(1 + r) at the end 1 (0) of the period and the expected return is kw = w1 µ1 + w2 r, w1 = nS V (0) , w2 = 1 − w1 . The assumption of normal joint returns is violated but the standard deviation of this portfolio can be easily computed directly from the

6. Product measures

185

definition giving σw = w1 σ1 (σ2 = 0 of course and the formula is consistent with the above).

Remark 6.3 The above considerations can be immediately generalized to portfolios built of any finite number of ingredients with the following key formulae X kw = wi k i , X µw = w i µi , X 2 σw = wi wj cij . i,j

This is just the beginning of the story started in the 1950s by Nobel prize winner Harry Markowitz. A vast number of papers and books on this topic have been written since, proving the general observation that ‘simple is beautiful’.

6.6 Proofs of propositions Proof (of Proposition 6.2) Denote by FR the σ-field generated by the Borel ‘rectangles’ R = {B1 × B2 : B1 , B2 ∈ B}, and by FI the σ-field generated by the true rectangles I = {I1 × I2 : I1 , I2 are intervals}. Since I ⊂ R, obviously FI ⊂ FR . To show the inverse inclusion we show that Borel cylinders B1 × Ω2 and Ω1 × B2 are in FI . For that write D = {A : A × Ω2 ∈ FI }, note that this is a σ-field containing all intervals hence B ⊂ D as required.

Proof (of Proposition 6.11) P Let sn = ck 1Ak be an increasing sequence of simple functions convergent to R f . Let Rk = Ak × [0, ck ] and the union of such rectangles is in fact sn dm. S S Then ∞ n=1 k Rk = Af so Af is measurable. For the second claim take a y section of Af which is the interval [0, f (x)). Its R measure is f (x) and by the definition of the product measure m2 (Af ) = f (x) dx.

186


Proof (of Proposition 6.16) The joint density is the product of the densities: fX,Y (x, y) = fX (x)fY (y) and substituting this to (6.6) immediately gives the result.

7 The Radon–Nikodym Theorem

In this chapter we shall consider the relationship between a real Borel measure ν and the Lebesgue measure m. Key to such relationships is Theorem 4.17, which shows that for each non-negative integrable real function f , the set function Z A 7→ ν(A) = f dm (7.1) A

defines a (Borel) measure ν on (R, M). The natural question to ask is the converse: exactly which real Borel measures can be found in this way? We shall find a complete answer to this question in this chapter, and in keeping with our approach in Chapters 5 and 6, we shall phrase our results in terms of general measures on an abstract set Ω.

7.1 Densities and Conditioning The results we shall develop in this chapter also allow us to study probability densities (introduced in Section 4.7.2), conditional probabilities and conditional expectations (see Sections 5.4.3 and 6.5.3) in much greater detail. For ν as R defined above to be a probability measure, we clearly require f dm = 1. In particular, if ν = PX is the distribution of a random variable X the function f = fX corresponding to ν in (7.1) was called the density of X. In similar fashion we defined the joint density f(X,Y ) of two random variables in Section 6.5.1, by reference of their joint distribution to two-dimensional 187

188


Lebesgue measure m2 : if X and Y are real random variables defined on some probability space (Ω, F, P ) their joint distribution is the measure defined on Borel subsets B of R2 by P(X,Y ) (B) = P ((X, Y ) ∈ B). In the special case where this measure, relative to m2 , is given as above by an integrable function f(X,Y ) , we say that X and Y have this function as their joint density. This, in turn, leads naturally (see Section 6.5.3) to the concepts of conditional density f(X,Y ) (a, y) h(y|a) = h(y|X = a) = fX (a) and conditional expectation Z E(Y |X = a) = yh(y|a) dy. R

RRecalling that X : Ω → R, the last equation can be written as E(Y |X)(ω) = yh(y|X(ω)) dy, displaying the conditional expectation as a random variable R E(Y |X) : Ω → R, measurable with respect to the σ-field FX generated by X. An application of Fubini’s theorem leads to a fundamental identity, valid for all A ∈ FX Z Z E(Y |X) dP = Y dP. (7.2) A

A

The existence of this random variable in the general case, irrespective of the existence of a joint density, is of great importance in both theory and applications – Williams [12] calls it ‘the central definition of modern probability’. It is essential for the concept of martingale, which plays such a crucial role in many applications, and which we introduce at the end of this chapter. As we described in Section 5.4.3, the existence of orthogonal projections in L2 allows one to extend the scope of the definition further still: instead of restricting ourselves to random variables measurable with respect to σ-fields of the form FX we specify any sub-σ-field G of F and ask for a G-measurable random variable E(Y |G) to play the role of E(Y |X) = E(Y |FX ) in (7.2). As was the case for product measures, the most natural context for establishing the properties of the conditional expectation is that of general measures; note that the proof of Theorem 4.17 simply required monotone convergence to establish the countable additivity of P. We therefore develop the comparison of abstract measures further, as always guided by the specific examples of random variables and distributions.

7.2 The Radon–Nikodym Theorem R In the special case where the measure ν has the form ν(A) = A f dm for some non-negative integrable function f we said (Section 4.7.2) that ν is absolutely

7. The Radon–Nikodym Theorem

189

R continuous with respect to m. It is immediate that A f dm = 0 whenever m(A) = 0 (see Theorem 4.3 (iv)). Hence m(A) = 0 implies ν(A) = 0 when the measure ν is given by a density. We use this as a definition for the general case of two given measures.

Definition 7.1 Let Ω be a set and let F be a σ-field of its subsets. (The pair (Ω, F) is a measurable space) Suppose that ν and µ are measures on (Ω, F). We say that ν is absolutely continuous with respect to µ if µ(A) = 0 implies ν(A) = 0 for A ∈ F . We write this as ν µ.

Exercise 7.1 Let λ1 , λ2 and µ be measures on (Ω, F). Show that if λ1 µ and λ2 µ then (λ1 + λ2 ) µ. It will not be immediately obvious what this definition has to do with the usual notion of continuity of functions. We shall see later in this chapter how it fits with the concept of absolute continuity of real functions. For the present, we note the following reformulation of the definition, which is not needed for the main result we will prove, but serves to bring the relationship between ν and µ a little ‘closer to home’ and is useful in many applications:

Proposition 7.1 Let ν and µ be finite measures on the measurable space (Ω, F). Then ν µ if and only if for every ε > 0 there exists a δ > 0 such that for F ∈ F, µ(F ) < δ implies ν(F ) < ε. Hint Suppose the (ε, δ)-condition fails. We can then find ε > 0 and sets (Fn ) such that for all n ≥ 1, µ(Fn ) < 21n but ν(Fn ) > ε. Consider µ(A) and T S ν(A) for A = n≥1 ( i≥n Fi ).

We generalise from the special case of Lebesgue measure: if R µ is any measure on (Ω, F) and f : Ω → R is a measurable function for which f dµ exists, then R ν(F ) = F f dµ defines a measure ν µ. (This follows exactly as for m, since R µ(F ) = 0 implies F f dµ = 0. Note that we employ the convention 0 × ∞ = 0.) For σ-finite measures, the following key result asserts the converse:

190


Theorem 7.2 (Radon–Nikodym) Given two σ-finite measures ν, µ on a measurable space (Ω, F), with ν µ, Rthen there is a non-negative measurable function h : Ω → R such that ν(F ) = h dµ for everyR F ∈ F. The function h is unique up to µ-null sets: if g also F satisfies ν(F ) = F g dµ for all F ∈ F, then g = h a.e. (µ).

Since the most interesting case for applications arises for probability spaces and then h ∈ L1 (µ), we shall initially restrict attention to the case where µ and ν are finite measures. In fact, it is helpful initially to take µ to be a probability measure, i.e. µ(Ω) = 1. From among several different approaches to this very important theorem, we base our argument on one given by R.C. Bradley in the American Mathematical Monthly (Vol 96, no 5., May 1989, pp. 437–440), since it offers the most ‘constructive’ and elementary treatment of which we are aware. It is instructive to begin with a special case: Suppose (until further notice) that µ(Ω) = 1. We say that the measure µ dominates ν when 0 ≤ ν(F ) ≤ µ(F ) for every F ∈ F. This obviously implies ν µ. In this simplified situation we shall construct the required function h explicitly. First we generalise the idea of partitions and their refinements, which we used to good effect in constructing the Riemann integral, to measurable subsets in (Ω, F).

Definition 7.2 Let (Ω, F) be a measurable space. A finite (measurable) partition of Ω is a finite collection of disjoint subsets P = (Ai )i≤n in F whose union is Ω. The finite partition P 0 is a refinementof P if each set in P is a disjoint union of sets in P 0 .

Exercise 7.2 Let P1 and P2 be finite partitions of Ω. Show that the coarsest partition (i.e. with least number of sets) which refines them both consists of all intersections A ∩ B, where A ∈ P1 , B ∈ P2 . The following is a simplified ‘Radon–Nikodym theorem’ for dominated measures:

Theorem 7.3 Suppose that µ(Ω) = 1 and 0 ≤ ν(F ) ≤ µ(F ) for every F ∈ F. Then there


191

exists a non-negative F-measurable function h on Ω such that ν(F ) = for all F ∈ F.

R

F

h dµ

We shall prove this in three steps: in Step 1 we define the required function hP for sets in a finite partition P and compare the functions hP1 and hP2 when R 2 the partition P2 refines P1 . This enables us to show that the integrals h dµ are non-decreasing if we Rtake successive refinements. Since they are Ω P also bounded above (by 1), c = sup Ω h2P dµ exists in R. In Step 2 we then construct the desired function h by a careful limit argument, using the convergence theorems of Chapter 4. In Step 3 we show that h has the desired properties. Step 1: The function hP for a finite partition Suppose that 0 ≤ ν(F ) ≤ µ(F ) for every F ∈ F. Let P = {A1 , A2 , . . . , Ak } be a finite partition of Ω such that each Ai ∈ F. Define the simple function hP : Ω → R by setting ν(Ai ) for ω ∈ Ai when µ(Ai ) > 0, and hP (ω) = 0 otherwise. µ(Ai ) R Since hP is constant on each ‘atom’ Ai , ν(Ai ) = Ai hP dµ. Then hP has the following properties: hP (ω) = ci =

(i) For each finite partition P of Ω, 0 ≤ hP (ω) ≤ 1 for all ω ∈ Ω. R S (ii) If A = j∈J A R j for an index set J ⊂ {1, 2, . . . k} then ν(A) = F hP dµ. Thus ν(Ω) = Ω hP dµ.

(iii) If P1 and P2 are finite partitions of Ω and P2 refines P1 then, with hn = hPn , (n = 1, 2) we have R R (a) for all A ∈ P1 , A h1 dµ = ν(A) = A h2 dµ, R R (b) for all A ∈ P1 , A h1 h2 dµ = A h21 dµ. R R (iv) Ω (h22 − h21 ) dµ = Ω (h2 − h1 )2 dµ and therefore Z

Ω

h22 dµ =

Z

Ω

h21 dµ +

Z

Ω

(h2 − h1 )2 dµ ≥

Z

Ω

h21 dµ.

We now prove these assertions in turn. (i) This is trivial by construction of hP , since µ dominates ν.

192


S (ii) Let A = j∈J Aj for some index set J ⊂ {1, 2, . . . , k}. Since the {Aj } are disjoint and ν(Aj ) = 0 whenever µ(Aj ) = 0, we have ν(A) =

X j∈J

X

=

j∈J,µ(Ai )>0

=

Z

A

ν(Aj ) µ(Aj ) µ(Aj ) j∈J,µ(Ai )>0 XZ cj µ(Aj ) = hP dµ X

ν(Aj ) =

j∈J

Aj

hP dµ.

In particular, since P partitions Ω, this holds for A = Ω. S (iii) (a) With the Pn , hn as above (n = 1, 2) we can write A = j∈J Bj for each A ∈ P1 , where J is a finite index set and Bj ∈ P2 . The sets Bj are pairwise disjoint, and again ν(Bj ) = 0 when µ(Bj ) = 0, so that Z X X ν(Bj ) µ(Bj ) h1 dµ = ν(A) = ν(Bj ) = µ(Bj ) A j∈J j∈J,µ(Bj )>0 Z XZ = h2 dµ = h2 dµ. Bj

j∈J

A

ν(A) (b) With A as in part (a) and µ(A) > 0, note that h1 = µ(A) is constant on A, so that Z Z Z Z ν(A) (ν(A))2 ν(A) 2 h1 h2 dµ = h2 dµ = = ( ) dµ = h21 dµ. µ(A) µ(A) µ(A) A A A A R (iv) By (iii) (b), A h1 (h2 − h1 ) dµ = 0 for every A ∈ P1 . Since the Ai ∈ P1 partition Ω, we also have

Z

Ω

h1 (h2 − h1 ) dµ =

k Z X i=1

Ai

h1 (h2 − h1 ) dµ = 0.

Hence Z

2

Ω

(h2 − h1 ) dµ = = =

Z

Z

Z

Ω

(h22 − 2h1 h2 + h21 ) dµ

Ω

[h22 − 2h1 (h2 − h1 ) − h21 ] dµ

Ω

(h22 − h21 ) dµ,

and thus Z

Ω

h22

dµ =

Z

Ω

h21

dµ +

Z

2

Ω

(h2 − h1 ) dµ ≥

Z

Ω

h21 dµ.


193

Step 2: Passage to the limit – construction of h. R 2 In Step 1 we showed that the integrals Ω hP dµ are non-decreasing over successive refinements of a finite partition of Ω. Moreover, by (i) above, each funcR tion hP satisfies 0 ≤ hP (ω) ≤ 1 for all ω ∈ Ω. Thus, setting c = sup Ω h2P dµ, where the supremum is taken over all finite partitions of Ω, we have 0 ≤ c ≤ 1. (Here we use the assumption that µ(Ω) = 1.) each n ≥ 1 let Pn be a finite measurable partition of Ω such that R For 2 h dµ > c− 41n . Let Qn be the smallest common refinement of the partitions Ω Pn P1 , P2 , . . . , Pn . For each n, Qn refines Pn by construction, and Qn+1 refines Qn since each Qk consists of all intersections A1 ∩ A2 ∩ . . . ∩ Ak , where Ai ∈ Pi , i ≤ k. Hence each set in Qn is a disjoint union of sets in Qn+1 . We therefore have the inequalities: Z Z Z 1 h2Pn dµ ≤ h2Qn dµ ≤ h2Qn+1 dµ ≤ c. c− n < 4 Ω Ω Ω Using the identity proved in Step 1 (iv), we now have Z Z 1 (hQn+1 − hQn )2 dµ = (h2Qn+1 − h2Qn ) dµ < n . 4 Ω Ω The Schwarz inequality applied with f = |hQn+1 − hQn | and g ≡ 1, then yields for each n ≥ 1, Z 1 |hQn+1 − hQn | dµ < n . 2 Ω R P By the Beppo–Levi Theorem, since n≥1 Ω |hQn+1 − hQn | dµ is finite, we P conclude that the series n≥1 (hQn+1 − hQn ) converges almost everywhere (µ), so that the limit function X (hQn+1 − hQn ) = lim hQn h = h P1 + n≥1

n

(noting that Q1 = P1 ) is well-defined almost everywhere (µ). We complete the construction by setting h = 0 on the exceptional µ-null set. Step 3: Verification of the properties of h. By Step 1 (i) it follows that 0 ≤ h(ω) ≤ 1, and it is clear from its construction that h is F-measurable. R We need to show that ν(F ) = F h dµ for every F ∈ F. Fix any such measurable set F and let n ≥ 1. Define Rn as the smallest common refinement of the two partitions Qn (defined as in Step 2)R and {F, F c }. Since F is a finite disjoint union of sets in Rn , we have ν(F ) = F hRn dµ from Step 1 (ii).

194


R R By Step 2,R c − 41n < Ω h2Qn dµ ≤ Ω h2Rn dµ ≤ c, so, as before, we can conclude that Ω (hRn − hQn )2 dµ < 41n , and using the Schwarz inequality once more, this time with g = 1F , we have Z Z 1 | (hRn − hQn ) dµ| ≤ |hRn − hQn | dµ < n . 2 F F R R R For all n, ν(F ) = F hRn dµ = F (hRn −hQn ) dµ+ F hQn dµ. The first integral R on the right converges to 0 as n → ∞, while the second converges to F h dµ by dominated convergence theorem (since for Rall n ≥ 1, 0 ≤ hQn ≤ 1 and µ(Ω) is finite). Thus we have verified that ν(F ) = F h dµ, as required.

It is straightforward to check that the assumption µ(Ω) = 1 is not essential since for any finite positive measure µ we can repeat the above arguments µ dν instead of µ. We write the function h defined above as dµ and call using µ(Ω) it the Radon–Nikodym derivative of ν with respect to µ. Its relationship to derivatives of functions will become clear when we consider real functions of bounded variation.

Exercise 7.3 Let Ω = [0, 1] with Lebesgue measure and consider measures µ, ν given by densities 1A , 1B respectively. Find a condition on the sets A, B so dν applying that µ dominates ν and find the Radon-Nikodym derivative dµ the above definition of the function h.

Exercise 7.4 Suppose Ω is a finite set equipped with the algebra of all subsets. Let µ and ν be two measures on Ω such that µ({ω}) 6= 0, ν({ω}) 6= 0, for all dν ω ∈ Ω. Decide under which conditions µ dominates ν and find dµ . The next observation is an easy application of the general procedure highlighted in Remark 4.1:

Proposition 7.4 dµ If µ and ϕ are finite measures with 0 ≤ µ ≤ ϕ, and if hµ = dϕ is constructed as above, then for any non-negative F-measurable function g on Ω we have Z Z g dµ = ghµ dϕ. Ω

Ω

The same identity holds for any g ∈ L1 (µ).


195

Hint Begin with indicator functions, use linearity of the integral to extend to simple functions, and monotone convergence for general non-negative g. The rest is obvious from the definitions. For finite measures we can now prove the general result announced earlier:

Theorem 7.5 (Radon–Nikodym) Let ν and µ be finite measures on the measurable space (Ω, F) and suppose that ν µ. Then there is a non-negative F-measurable function h on Ω such R that ν(A) = A h dµ for all A ∈ F.

Proof Let ϕ = ν + µ. Then ϕ is a positive finite measure which dominates both dµ dν and hµ = dϕ are ν and µ. Hence the Radon–Nikodym derivatives hv = dϕ well-defined by the earlier constructions. Consider the sets F = {hµ > 0} and R G = {hµ = 0} in F. Clearly µ(G) = G hµ dϕ = 0, hence also ν(G) = 0, since ν µ. Define h = hhµν 1F , and let A ∈ F, A ⊂ F . By the previous proposition, with h1A instead of g, we have Z Z Z ν(A) = hν dϕ = hhµ dϕ = h dµ A

A

A

as required. Since µ and ν are both null on G this proves the theorem.

Exercise 7.5 Let Ω = [0, 1] with Lebesgue measure and consider probability measures µ, ν given by densities f, g respectively. Find a condition characterising the absolute continuity ν µ and find the Radon-Nikodym derivative dν dµ .

Exercise 7.6 Suppose Ω is a finite set equipped with the algebra of all subsets and let µ and ν be two measures on Ω. Characterise the absolute continuity dν ν µ and find dµ . You can now easily complete the picture for σ-finite measures and verify that the function h is ‘essentially unique’:

196


Proposition 7.6 The Radon–Nikodym theorem remains valid if the measures ν and µ are σfinite: for any two such measures with ν µ we can find R a finite-valued nonnegative measurable function f on Ω such that ν(F ) = F h dµ for all F ∈ F. + The function h so R defined is unique up to µ-null sets, i.e. if g : Ω → R also satisfies ν(F ) = Ω g dµ for all F ∈ F then g = h a.e. (with respect to µ).

Hint There are sequences (An ), (Bm ) of sets in F with µ(An ), ν(Bm ) finite S S for all m, n ≥ 1 and n≥1 An = Ω = m≥1 Bm . We can choose these to be sequences of disjoint sets (why?). Hence display Ω as the disjoint union of the sets An ∩ Bm (m, n ≥ 1), thus finding a sequence (Cn ) of disjoint sets with union Ω, all of whose members have finite measure under both µ and ν. Fix n and apply the above results to the measurable space (Ω, Fn ), where Fn = {F ∩ Cn : F ∈ F}, then ‘paste together’ the resulting functions for all n. Radon-Nikodym derivatives of measures obey simple combination rules which follow from the uniqueness property. We illustrate this with the sum and composition of two Radon-Nikodym derivatives, and leave the ’inverse rule’ as an exercise.

Proposition 7.7 Assume we are given σ-finite measures λ, ν, µ satisfying λ µ and ν µ with dλ dν Radon-Nikodym derivatives dµ and dµ , respectively. (i) With φ = λ + ν we have (ii) If λ ν then

dλ dµ

=

dλ dν dν dµ

dφ dµ

=

dλ dµ

+

dν dµ

a.s. (µ),

a.s. (µ).

Exercise 7.7 Show that if µ, ν are equivalent measures, i.e. both ν µ and µ ν are true, then dµ dν = ( )−1 a.s. (µ). dν dµ Given a pair of σ-finite measures λ, µ on (Ω, F) it is natural to ask whether we can identify the sets for which µ(E) = 0 implies λ(E) = 0. This would mean that we can split the mass of λ into two pieces, one being represented by a µ-integral, and the other ‘concentrated’ on µ-null sets, i.e. away from the mass of µ. We turn this idea of ‘separating’ the masses of two measures into the following


197

Definition 7.3 If there is a set E ∈ F such that λ(F ) = λ(E ∩ F ) for every F ∈ F then λ is concentrated on E. If two measures µ, ν are concentrated on disjoint subsets of Ω, we say that they are mutually singular and write µ ⊥ ν. Clearly, if λ is concentrated on E and E∩F = Ø, then λ(F ) = λ(E∩F ) = 0. Conversely, if for all F ∈ F, F ∩ E = Ø implies λ(F ) = 0, consider λ(F ) = λ(F ∩ E) + λ(F E). Since (F E) ∩ E = Ø we must have λ(F E) = 0, so λ(F ) = λ(F ∩ E). We have proved that λ is concentrated on E if and only if for all F ∈ F, F ∩ E = Ø implies λ(F ) = 0. We gather some simple facts about mutually singular measures:

Proposition 7.8 If µ, ν, λ1, λ2 are measures on a σ-field F, the following are true:

(i) If λ1 ⊥ µ and λ2 ⊥ µ then also (λ1 + λ2 ) ⊥ µ.

(ii) If λ1 µ and λ2 ⊥ µ then λ2 ⊥ λ1 .

(iii) If ν µ and ν ⊥ µ then ν = 0.

Hint For (i), with i = 1, 2 let Ai , Bi be disjoint sets with λi concentrated on Ai , µ on Bi . Consider A1 ∪ A2 and B1 ∩ B2 . For (ii) use the remark preceding the proposition. The next result shows that a unique ‘mass splitting’ of a σ-finite measure relative to another is always possible:

Theorem 7.9 (Lebesgue decomposition) Let λ, µ be σ-finite measures on (Ω, F). Then λ can be expressed uniquely as a sum of two measures, λ = λa + λs where λa µ and λs ⊥µ.

Proof Existence: We consider finite measures; the extension to the σ-finite case is routine. Since 0 ≤ λ ≤ λ + µ = ϕ, i.e. φ dominates λ, there is 0 ≤ h ≤ 1 such R that λ(E) = E h dϕ for all measurable E. Let A = {ω : h(ω < 1} and B = {ω : h(ω) = 1}. Set λa (E) = λ(A ∩ E) and λs (E) =R λ(B ∩ E) for R every E ∈ F. h dλ, so that Now if E ⊂ A and µ(E) = 0 then λ(E) = h dϕ = E E R (1 − h) dλ = 0. But h < 1 on A, hence also on E. Therefore we must have E λ(E) = 0. Hence if E ∈ F and µ(E) = 0, λa (E) = λ(A ∩ E) = 0 as A ∩ E ⊂ A.

198


R So = E h dϕ = R λa λ. On the other hand, if E ⊂ B we obtain λ(E) c E 1 d(λ + µ) = λ(E) + µ(E), so that µ(E) = 0. As A = B we have shown that µ(E) = 0 whenever E ∩ A = Ø, so that µ is concentrated on A. Since λs is concentrated on B this shows that λs and µ are mutually singular. Uniqueness is left to the reader. (Hint: employ Proposition 7.8.) The theorem is proved. Combining this with the Radon–Nikodym theorem we can describe the structure of λ with respect to µ as ‘basis measure’:

Corollary 7.10 With µ, λ, λa , λs as in the theorem, there is a µ-a.s. unique non-negative meaR surable function h such that λ(E) = E h dµ + λs (E) for every E ∈ F.

Remark 7.1 This result is reminiscent of the structure theory of finite-dimensional vector Pm spaces: if x ∈ Rn and m < n, we can write x = y + z, where y = i=1 yi ei is the orthogonal projection onto Rm and z is orthogonal to this subspace. We also exploited similar ideas for Hilbert space. In this sense the measure µ has the role of a ‘basis’ providing the ‘linear combination’ which describes the projection of the measure λ onto a subspace of the space of measures on Ω.

Exercise 7.8 Consider the following measures on the real line: P1 = δ0 , P2 = 1 1 1 25 m|[0,25] , P3 = 2 P1 + 2 P2 (see Example 3.1). For which i 6= j do we have Pi Pj ? Find the Radon-Nikodym derivative in each such case.

Exercise 7.9 Let λ = δ0 + m|[1,3] , µ = δ1 + m|[2,4] and find λa , λs , and h as in Corollary 7.10.

7.3 Lebesgue–Stieltjes measures Recall (see Section 3.5.3) that given any random variable X : Ω −→ R, we define its probability distribution as the measure PX = P ◦ X −1 on Borel


199

sets on R (i.e. we set P (X ≤ x) = P ◦ X −1 ((−∞, x]) = PX ((−∞, x]) and extend this to B.) Setting FX (x) = PX ((−∞, x]) we verified in Proposition 4.30 that the distribution function FX so defined is monotone increasing, right-continuous, with limits at infinity FX (−∞) = limx→−∞ FX (x) = 0 and FX (+∞) = limx→∞ FX (x) = 1. R x In Chapter 4 we studied the special case where FX (x) = PX ((−∞, x]) = f dm for some real function fX , the density of PX with respect to −∞ X Lebesgue measure m, Proposition 4.22 showed that if fX is continuous, then FX is differentiable and has the density fX as its derivative at every x ∈ R. On the other hand, the Lebesgue function in Example 4.8 illustrated that continuity of FX is not sufficient to guarantee the existence of a density. Moreover, when FX has a density fX , the measure PX was said to be ‘absolutely continuous’ with respect to m. In the context of the Radon–Nikodym theorem we should reconcile the terminology of this special case withR the general one considered in the present Chapter. Trivially, when PX (B) = B fX dm X we have PX m, so that PX has a Radon–Nikodym derivative dP dm with dPX respect to m. The a.s. uniqueness ensures that dm = fX a.s. Later in this chapter we shall establish the precise analytical requirements on the cumulative distribution function FX which will guarantee the existence of a density.

7.3.1 Construction of Lebesgue–Stieltjes measures To do this we first study, only slightly more generally, measures defined on (R, B) which correspond in similar fashion to increasing, right-continuous functions on R. Their construction mirrors that of Lebesgue measure, with only a few changes, by generalising the concept of ‘interval length’. The measures we obtain are known as Lebesgue-Stieltjes measures. In this context we call a function F : R −→ R a distribution function if F is monotone increasing and right-continuous. It is clear that every finite measure µ defined on (R, B) defines such a function by F (x) = µ((−∞, x]), with F (−∞) = limx→−∞ F (x) = 0, F (+∞) = limx→∞ F (x) = µ(Ω). Our principal concern, however, is with the converse: given a monotone right-continuous F : R −→ R, can we always associate with F a measure on (Ω, B), and if so, what is its relation to Lebesgue measure? The first question is answered by looking back carefully at the construction of Lebesgue measure m on R in Chapter 2: first we defined the natural concept of interval length, l(I) = b − a, for any interval I with endpoints a, b (a < b), and by analogy with our discussion of null sets, we defined Lebesgue outer measure m∗ for an arbitrary subset of R as the infimum of the total lengths

200


P∞

n=1 l(In ) of all sequences (In )n≥1 of intervals covering A. To generalise this idea, we should clearly replace b − a by F (b) − F (a) to obtain a ’generalised interval length’ relative to F , but since F is only right-continuous we will need to take care of possible discontinuities. Thus we need to identify the possible discontinuities of monotone increasing functions – fortunately these functions are rather well-behaved, as you can easily verify in the following:

Proposition 7.11 If F : R −→ R is monotone increasing (i.e. x1 ≤ x2 implies F (x1 ) ≤ F (x2 )) then the left-limit F (x−) and the right-limit F (x+) exist at every x ∈ R and F (x−) ≤ F (x) ≤ F (x+). Hence F has at most countably many discontinuities, and these are jump discontinuities, i.e. F (x−) < F (x+). Hint For any x, consider sup{F (y) : y < x} and inf{F (y) : x < y} to verify the first claim. For the second, note that F (x−) < F (x+) if F has a discontinuity at x. Use the fact that Q is dense in R to show that there can only be countably many such points. Since F is monotone, it remains bounded on bounded sets. For simplicity we assume that limx→−∞ F (x) = 0. We define the ‘length relative to F ’ of the bounded interval (a, b] by lF (a, b] = F (b) − F (a). Note that we have restricted ourselves to left-open, right-closed intervals. Since F is right-continuous, F (x+) = F (x) for all x, including a, b. Thus lF (a, b] = F (b+) − F (a+), and all jumps of F have the form F (x) − F (x−). By restricting to intervals of this type we also ensure that lF is additive over adjoining intervals: if a < c < b then lF (a, b] = lF (a, c] + lF (c, b]. We generalise Definition 2.2 as follows:

Definition 7.4 The F -outer measure of any set A ⊆ R is the element of [0, ∞] m∗F (A) = inf ZF (A) where ZF (A) = {

∞ X

n=1

lF (In ) : In = (an , bn ], an ≤ bn , A ⊆

∞ [

n=1

In }.


201

Our ‘covering intervals’ are now also restricted to be left-open and rightclosed. This is essential to ‘make things fit together’, but does not affect measurability: recall (Theorem 2.16) that the Borel σ-field is generated whether we start from the family of all intervals or from various sub-families. Now consider the proof of Theorem 2.4 in more detail: our purpose there was to prove that the outer measure of an interval equals its length. We show how to adapt the proof to make this claim valid for m∗F and lF applied to intervals of the form (a, b]. It will be therefore helpful to review the proof of Theorem 2.4 before reading on! Step 1. The proof that m∗F ((a, b]) ≤ lF (a, b] remains much the same: To see that lF (a, b] ∈ ZF ((a, b]), we cover (a, b] by (In ) with I1 = (a, b], In = (a, a] = Ø, n > 1. The total length of this sequence is F (b) − F (a) = lF (a, b], hence the result follows by definition of inf. Step 2. It remains to show that lF (a, b] ≤ m∗F ((a, b]). Here we need to be careful always to ‘approach points from the right’ in order to make use of the right-continuity of F and thus to avoid its jumps. Fix ε > 0 and 0 < δ < b − a. By definition of inf we can find a covering of P∞ I = (a, b] by intervals In = (an , bn ] such that n=1 lF (In ) < m∗F (I) + 2ε . Next, let Jn = (an , b0n ), where by right-continuity of F , for each n ≥ 1 we can choose ε b0n > bn and F (b0n ) − F (bn ) < 2n+1 . Then F (b0n ) − F (an ) < {F (bn ) − F (an )} + ε 2n+1 . The (Jn )n≥1 then form an open cover of the compact interval [a + δ, b], so that by the Heine–Borel Theorem there is a finite subfamily (Jn )n≤N , which also covers [a + δ, b]. Re-ordering these N intervals Jn we can assume that their right-hand endpoints form an increasing sequence and then F (b) − F (a + δ) = lF (a + δ, b] ≤
0, hence F (b) − F (a + δ) ≤ m∗F (I) for every δ > 0. By right-continuity of F , letting δ ↓ 0 we obtain lF (a, b] = limδ↓0 lF (a + δ, b] ≤ m∗F (a, b]. This completes the proof that m∗F ((a, b]) = lF (a, b]. This is the only substantive change needed from the construction that led to Lebesgue measure. The proof that m∗F is an outer measure, i.e. m∗F (A) ≥ 0,

m∗F (Ø) = 0,

m∗F (A) ≤ m∗F (B) if A ⊆ B,

202


m∗F (

∞ [

i=1

Ai ) ≤

∞ X

m∗F (Ai ),

i=1

is word-for-word identical with that given for Lebesgue outer measure (Proposition 2.3, Theorem 2.5). Hence, as in Definition 2.3 we say that a set E is measurable (for the outer measure m∗F ) if for each A ⊆ R m∗F (A) = m∗F (A ∩ E) + m∗F (A ∩ E c ). Again, the proof of Theorem 2.8 goes through verbatim, and we denote the resulting Lebesgue–Stieltjes measure, i.e. m∗F restricted to the σ-field MF of Lebesgue–Stieltjes measurable sets, by mF . By construction, just like Lebesgue measure, mF is a complete measure: subsets of mF -null sets are in MF . However, as we shall see later, MF does not always coincide with the σ-field M of Lebesgue-measurable sets, although both contain all the Borel sets. It is also straightforward to verify that the properties of Lebesgue measure proved in Section 2.4 hold for general Lebesgue–Stieltjes measures, with one exception: the outer measure m∗F will not, in general, be translation-invariant. We can see this at once for intervals, since lF ((a + t, b + t]) = F (b + t) − F (a + t) will not usually equal F (b) − F (a); simply take F (x) = x3 , for example. In fact, it can be shown that Lebesgue measure is the unique translation-invariant measure on R. Note, moreover, that a singleton {a} is now not necessarily a null set for mF : we have, by the analogue of Theorem 2.13, that mF ({a}) = lim mF ((a − n→∞

1 1 , a]) = F (a) − lim F (a − ) = F (a) − F (a−). n→∞ n n

Thus, the measure of the set {a} is precisely the size of the jump at a (if any). From this it is easy to see by similar arguments how the ‘length’ of an interval depends on the presence or absence of its endpoints: given that mF ((a, b]) = F (b)−F (a), we see that: mF ((a, b)) = F (b−)−F (a), mF ([a, b]) = F (b)−F (a−), mF ([a, b)) = F (b−) − F (a).

Example 7.1 When F = 1[a,∞) we obtain mF = δa , the Dirac measure concentrated at a. Similarly, we can describe a general discrete probability distribution, where the random variable X takes the values {ai : i = 1, 2, . . . , n} with probabilities {pi = 1, 2, . . . , n} as the Lebesgue–Stieltjes measure arising from the function Pn F = i=1 pi 1[ai ,∞) .

Mixtures of discrete and continuous distributions, such as described in Example 3.1, clearly also fit into this picture. Of course, Lebesgue measure m is


203

the special case where the distribution is uniform, i.e. if F (x) = x for all x ∈ R then mF = m.

Example 7.2 Only slightly more generally, every finite Borel measure µ on R corresponds to a Lebesgue–Stieltjes measure, since the distribution function F (x) = µ((−∞, x]) is obviously increasing and is right-continuous by Theorem 2.13 applied to µ and the intervals In = (−∞, x + n1 ]). The corresponding Lebesgue–Stieltjes measure mF = µ, since they coincide on the generating family of intervals of the above form. Hence they coincide on the σ-field B of Borel sets. By our construction of mF as a complete measure it follows that mF is the completion of µ.

Example 7.3 Return to the Lebesgue function F discussed in Example 4.8. Since F is continuous and monotone increasing, it induces a Lebesgue–Stieltjes measure mF on the interval [0, 1], whose properties we now examine. On each ‘middle thirds’ set F is constant, hence these intervals are null sets for mF , and as there are countably many of them, so is their union, the ‘middle thirds’ set, D. Hence the Cantor set C = D c satisfies 1 = F (1) − F (0) = mF ([0, 1]) = mF (C) (Note that since F is continuous, mF ({0}) = F (0) − F (0−) = 0; in fact, each singleton is mF -null.) We thus conclude that mF is concentrated on a null set for Lebesgue measure m, i.e. mF ⊥ m, and that in the Lebesgue decomposition of mF relative to m there is no absolutely continuous component (by uniqueness of the decomposition).

Exercise 7.10 Suppose the monotone increasing function F is non-constant at most countably many points (as would be the case for a discrete distribution). Show that every subset of R is mF -measurable. Hint Consider mF over the bounded interval [−M, M ] first.

204


Exercise 7.11 Find the Lebesgue-Stieltjes measure mF generated by   0 if x < 0, F (x) = 2x if x ∈ [0, 1],  2 if x ≥ 1.

7.3.2 Absolute continuity of functions We now address the requirements on a distribution F which ensure that it has a density. As we saw in Example 4.8, continuity of a probability distribution function does not guarantee the existence of a density. The following stronger restriction, however, does the trick:

Definition 7.5 A real function F is absolutely continuous on the interval [a, b] if, given ε > 0, there is δ > 0 such that for every finite set of disjoint intervals Jk = (xk , yk ), k ≤ P P n, contained in [a, b] and with nk=1 (yk − xk ) < δ, we have nk=1 |F (xk ) − F (yk )| < ε. This condition will allow us to identify those distribution functions which generate Lebesgue–Stieltjes measures that are absolutely continuous (in the sense of measures) relative to Lebesgue measure. We will see shortly that absolutely continuous functions are also ‘of bounded variation’: this describes functions which do not ‘vary too much’ over small intervals. First we verify that the indefinite integral (see Proposition 4.22) relative to a density is absolutely continuous.

Proposition 7.12 1 RIfxf ∈ L ([a, b]), where the interval [a, b] is finite, then the function F (x) = a f dm is absolutely continuous.

Hint Use the absolute continuity of µ(G) = measure m.

R

G

|f | dm with respect to Lebesgue

Exercise 7.12 Decide which of the following functions are absolutely continuous: (a)


f (x) = |x|, x ∈ [−1, 1], (b) g(x) = function.

205

√

x, x ∈ [0, 1], (c) the Lebesgue

The next result is the important converse to the above example, and shows that all Stieltjes integrals arising from absolutely continuous functions lead to measures which are absolutely continuous relative to Lebesgue measure, and hence have a density. Together with the Example this characterises the distributions arising from densities (under the conditions we have imposed on distribution functions).

Theorem 7.13 If F is monotone increasing and absolutely continuous on R, let mF be the Lebesgue–Stieltjes measure it generates. Then every Lebesgue-measurable set is mF -measurable, and on these sets mF m.

Proof We first show that if the Borel set B has m(B) = 0, then also mF (B) = 0. Recall that, given δ > 0 we can find an open set O containing B with m(O) < δ (Theorem 2.12), and there is a sequence of disjoint open intervals (Ik )k≥1 , Ik = (ak , bk ) with union O. Since the intervals are disjoint, their total length is less than δ. By the absolute continuity of F , given any ε > 0, we can find δ > 0 such that for every finite sequence of intervals Jk = (xk , yk ), k ≤ n, with total P P length nk=1 (yk − xk ) < δ, we have nk=1 {F (yk ) − F (xk )} < ε2 . Applying this P to the sequence (Ik )k≤n for a fixed n we obtain nk=1 {F (bk ) − F (ak )} < 2ε . P∞ As this holds for every n, we also have k=1 {F (bk ) − F (ak )} ≤ 2ε < ε. This is the total length of a sequence of disjoint intervals covering O ⊃ B, hence mF (B) < ε for every ε > 0, so mF (B) = 0. Now for every Lebesgue-measurable set E with m(E) = 0 we can find a Borel set B ⊇ E with m(B) = 0. Thus also mF (B) = 0. Now E is a subset of an mF -null set, hence it is also mF -null. Hence all m-measurable sets are mF measurable and m-null sets are mF -null, i.e. mF m when both are regarded as measures on M. Together with the Radon-Nikodym Theorem, the above result helps to clarify the structural relationship between Lebesgue measure and Lebesgue– Stieltjes measures generated by monotone increasing right-continuous functions, and thus, in particular, for probability distributions: when the function F is absolutely continuous it has a density f , and can therefore be written

206


as its ‘indefinite integral’. Since the Lebesgue–Stieltjes measure mF m, the F Radon–Nikodym derivative dm dm is well-defined. Conversely, for the density f of F to exist, the function F must be absolutely continuous. It now remains to F clarify the relationship between the Radon–Nikodym derivative dm dm and the density f . It is natural to expect from the example of a continuous f (Proposition 4.22) that f should be the derivative of F (at least m-a.e.). So we need to understand which conditions on F will ensure that F 0 (x) exists for m-almost all x ∈ R. We shall address this question in the somewhat wider context where the ‘integrator’ function F is no longer necessarily monotone increasing, but has bounded variation, as introduced in the next section.

7.3.3 Functions of bounded variation Since in general we need to handle set functions that can take negative values, for example, the map Z E −→ g dm, where g ∈ L1 (m), E

we therefore need a concept of ‘generalised length functions’ which are expressed as the difference of two monotone increasing functions. We need first to characterise such functions. This is done by introducing the following

Definition 7.6 A real function F is of bounded variation on [a, b] (briefly F ∈ BV [a, b]) if TF [a, b] < ∞, where for any x ∈ [a, b] TF [a, x] = sup{

n X k=1

|F (xk ) − F (xk−1 )|}

with the supremum taken over all finite partitions of [a, x] with a = x0 < x1 < . . . < xn = x. We introduce two further non-negative functions by setting PF [a, x] = sup{

n X k=1

and NF [a, x] = sup{

[F (xk ) − F (xk−1 )]+ }

n X k=1

[F (xk ) − F (xk−1 )]− }


207

where the supremum is again taken over all partitions of [a, x]. The functions TF (PF , NF ) are known respectively as the total (positive, negative) variation functions of F . We shall keep a fixed in what follows, and consider these as functions of x for x ≥ a. We can easily verify the following basic relationships between these definitions:

Proposition 7.14 If F is of bounded variation on [a, b], we have F (x) − F (a) = PF (x) − NF (x), while TF (x) = PF (x) + NF (x) for x ∈ [a, b]. Pn Pn Hint Consider p(x) = k=1 [F (xk ) − F (xk−1 )]+ and n(x) = k=1 [F (xk ) − F (xk−1 )]− for a fixed partition of [a, x] and note that F (x) − F (a) = p(x) − n(x). Now use the definition of the supremum. For the second identity consider TF (x) ≥ p(x) + n(x) = 2p(x) − F (x) + F (a) and use the first identity.

Proposition 7.15 If F is of bounded variation and a ≤ x ≤ b then TF [a, b] = TF [a, x] + TF [x, b]. Similar results hold for PF and NF . Hence all three variation functions are monotone increasing in x for fixed a ∈ R. Moreover, if F has bounded variation on [a, b], then it has bounded variation on any [c, d] ⊂ [a, b]. Hint Adding a point to a partition will increase all three sums. On the other hand, putting together partitions of [a, c] and [c, b] we obtain a partition of [a, b]. We show that bounded variation functions on finite intervals are exactly what we are looking for:

Theorem 7.16 Let [a, b] be a finite interval. A real function is of bounded variation on [a, b] if and only if it is the difference of two monotone increasing real functions on [a, b].

Proof If F is of bounded variation, use F (x) = [F (a) + PF (x)] − NF (x) from Proposition 7.14 to represent F as the difference of two monotone increasing functions.

208


Conversely, if F = g − h is the difference of two monotone increasing functions, then for any partition a = x0 < x1 < . . . < xn = b of [a, b] we obtain, since g, h are increasing, n X i=1

|F (xi ) − F (xi−1 )| = ≤

n X

i=1 n X i=1

|g(xi ) − h(xi ) − g(xi−1 ) + h(xi−1 )| [g(xi ) − g(xi−1 )] +

n X i=1

[h(xi ) − h(xi−1 )]

≤ g(b) − g(a) + h(b) − h(a).

Thus M = g(b) − g(a) + h(b) − h(a) is an upper bound independent of the choice of partition, and so TF [a, b] ≤ M < ∞, as required. This decomposition is minimal: if F = F1 − F2 and F1 , F2 are increasing, then for any partition a = x0 < x1 . . . < xn = b we can write, for fixed i ≤ n

{F (xi ) − F (xi−1 )}+ − {F (xi ) − F (xi−1 )}− = F (xi ) − F (xi−1 ) = {F1 (xi ) − F1 (xi−1 )} − {F2 (xi ) − F2 (xi−1 )}

which shows from the minimality property of x = x+ − x− that each term in the difference on the right dominates its counterpart of the left. Adding and taking suprema we conclude that PF is dominated by the total variation of F1 and NF by that of F2 . In other words, in the collection of increasing functions whose difference is F , the functions (F (a) + PF ) and NF have the smallest sum at every point of [a, b].

Exercise 7.13 (a) Let F be monotone increasing on [a, b]. Find TF [a, b]. (b) Prove that if F ∈ BV [a, b] then F is continuous a.e. (m) and Lebesgue-measurable. (c) Find a differentiable function which is not in BV [0, 1]. (d) Show that if there is a (Lipschitz) constant M > 0 such that |F (x) − F (y)| ≤ M |x − y| for all x, y ∈ [a, b], then F ∈ BV [a, b]. The following simple facts link bounded variation and absolute continuity for functions on a bounded interval [a, b]:

Proposition 7.17 Suppose the real function F is absolutely continuous on [a, b]; then we have:


209

(i) F ∈ BV [a, b],

(ii) If F = F1 − F2 is the minimal decomposition of F as the difference of two monotone increasing functions described in Theorem 7.16, then both F1 and F2 are absolutely continuous on [a, b]. Hint Given ε > 0 choose δ > 0 as in Definition 7.5. In (i), starting with an arbitrary partition (xi ) of [a, b] we cannot use the absolute continuity of F unless we know that the subintervals are of length δ. So add enough new partition points to guarantee this and consider the sums they generate. For (ii), compare the various variation functions when summing over a partition where the sum of intervals lengths is bounded by δ.

Definition 7.7 If F ∈ BV [a, b], where a, b ∈ R, let F = F1 − F2 be its minimal decomposition into monotone increasing functions. Define the Lebesgue–Stieltjes signed measure of F as the countably additive set function mF given on the σ-field B of Borel sets by mF = mF1 − mF2 , where mFi is the Lebesgue–Stieltjes measure of Fi , (i = 1, 2). We shall examine signed measures more generally in the next section. For the present, we note the following

Example 7.4 R When considering the measure PX (E) = E fX dm induced on R by a density fX we restrict attention to fX ≥ 0 to ensure that PX is non-negative. But for a measurable function f : R −→ R we set (Definition 4.4) Z Z Z Z + − f dm = f dm − f dm whenever |f | dm < ∞. E

E

E

E

R

The set function ν defined by ν(E) = E f dm then splits naturally R + into the + − + differenceRof two measures, i.e. ν = ν − ν , where ν (E) = E f dm and ν − (E) = R E f − dm. Restricting to a function f supported on [a, b] and setting x F (x) = a f dm we obtain mF = ν, and if F = F1 − F2 as in the above definition, then mF1 = ν + , mF2 = ν − by the minimality properties of the splitting of F .

210


7.3.4 Signed measures The above example and the definition of Lebesgue–Stieltjes measures generated by a BV function motivate the following abstract definition and the subsequent search for a similar decomposition into the difference of two measures. We proceed to outline briefly the structure of signed measures in the abstract setting, which provides a general context for the above development of Stieltjes integrals and distribution functions. Our results will enable us to define integrals of functions relative to signed measures by reference to the decomposition of the signed measure into ‘positive and negative parts’, exactly as above. We also obtain a more general Lebesgue decomposition and Radon–Nikodym theorem, thus completing the description of the structure of a bounded signed measure relative to a given σ-finite measure. This leads to the general version of the Fundamental Theorem of the Calculus signalled earlier.

Definition 7.8 A signed measure on a measurable space (Ω, F) is a set function ν : F −→ (−∞, +∞] satisfying (i) ν(Ø) = 0 S∞ P∞ (ii) ν( i=1 Ei ) = i=1 ν(Ei ) if Ei ∈ F and Ei ∩ Ej = Ø for i 6= j.

We need to avoid ambiguities like ∞−∞ by demanding that ν should take at most one of the values ±∞; therefore we consistently demand that ν(E) > −∞ for all sets E in its domain. Note also that in (ii) either both sides are +∞, or they are both finite, so that the series converges in R. Since the left side is unaffected by any re-arrangement of the terms of the series, it follows that the P∞ series converges absolutely whenever it converges, i.e. i=1 |ν(Ei )| < ∞ if and S∞ only if |ν( i=1 Ei )| < ∞. The convergence is clear in the motivating example, since for any E ⊆ R we have Z Z |ν(E)| = | f dm| ≤ |f | dm < ∞ when f ∈ L1 (R). E

E

Note that ν is finitely additive (let Ei = Ø for all i > n in (ii), then (i) implies Sn Pn ν( i=1 Ei ) = i=1 ν(Ei ) if Ei ∈ F and Ei ∩ Ej = Ø for i 6= j, i, j ≤ n). Hence if F ⊆ E, F ∈ F, and |ν(E)| < ∞, then |ν(F )| < ∞, since both sides of ν(E) = ν(F ) + ν(E \ F ) are finite and ν(EF ) > −∞ by hypothesis. Signed measures do not inherit the properties of measures without change: as a negative result we have


211

Proposition 7.18 A signed measure ν defined on a σ-field F is monotone increasing (F ⊂ E implies ν(F ) ≤ ν(E)) if and only ν is a measure on F. Hint Ø is a subset of every E ∈ F ! On the other hand, a signed measure attains its bounds at some sets in F. More precisely: given a signed measure ν on (Ω, F) one can find sets A and B in F such that ν(A) = inf{ν(F ) : F ∈ F} and ν(B) = sup{ν(F ) : F ∈ F}. Rather than prove this result directly we shall deduce it from the Hahn– Jordan decomposition theorem. This basic result shows how the set A and its complement can be used to define two (positive) measures ν + , ν − such that ν = ν + − ν − , with ν + (F ) = ν(F ∩ Ac ) and ν − (F ) = −ν(F ∩ A) for all F ∈ F. The decomposition is minimal: if ν = λ1 − λ2 where the λi are measures, then ν + ≤ λ1 and ν − ≤ λ2 . Restricting attention to bounded signed measures (which suffices for applications to probability theory), we can derive this decomposition by applying the Radon–Nikodym theorem. (Our account is a special case of the treatment given in [10], Ch.6, for complex-valued set functions.) First, given a bounded signed measure ν : F −→ R, we seek the smallest (positive) measure µ that dominates ν, i.e. satisfies µ(E) ≥ |ν(E)| for all E ∈ F. Defining |ν|(E) = sup{

∞ X i=1

|ν(Ei )| : {Ei } ⊂ F, E =

[

i≥1

Ei , Ei ∩ Ej = Ø if i 6= j}

produces a set function which satisfies |ν|(E) ≥ |ν(E)| for every E. The requirement µ(Ei ) ≥ |ν(Ei )| for all i then yields µ(E) =

∞ X i=1

µ(Ei ) ≥

∞ X i=1

|ν(Ei )|

for any measure µ dominating ν. Hence to prove that |ν| has the desired properties we only need to show that it is countably additive. We call |ν| the total variation of ν. Note that we use countable partitions of Ω here, just as we used sequences of intervals when defining Lebesgue measure in R.

Theorem 7.19 The total variation |ν| of a bounded signed measure is a (positive) measure on F.

212


Proof Partitioning E ∈ F into sets {Ei }, choose (ai ) in R+ such that ai < |ν|(Ei ) for all i. Partition each Ei in turn into sets {Aij }j , and by definition of sup P we can choose these to ensure that ai < |ν(Aij )| for every i ≥ 1. But P Pj the {Aij } also partition E, hence i ai < i,j |ν(Ai,j )| < |ν|(E). Taking the supremum over all sequences (ai ) satisfying these requirements ensures that P P i |ν|(Ei ) = sup i ai ≤ |ν|(E). For the converse inequality consider any partition {Bk } of E and note that for fixed k, {Bk ∩Ei }i≥1 partitions Bk , while for fixed i, {Bk ∩Ei }k≥1 partitions Ei . This means that X X X XX |ν(Bk )| = | ν(Bk ∩ Ei )| ≤ |ν(Bk ∩ Ei )|. k≥1

k≥1 i≥1

k≥1 i≥1

Since the terms of the double series are all non-negative, we can exchange the order of summation, so that finally X XX X |ν(Bk )| ≤ |ν(Bk ∩ Ei )| ≤ |ν|(Ei ). i≥1 k≥1

k≥1

i≥1

But the partition {Bk } of E was arbitrary, so the estimate on the right also dominates |ν|(E). This completes the proof that |ν| is a measure. We now define the positive (resp. negative) variation of the signed measure ν by setting: 1 (|ν| + ν), ν − = 2 Clearly both are positive measures on F, and ν+ =

1 (|ν| − ν). 2 we have

ν = ν + − ν − and |ν| = ν + + ν − . With these definitions we can immediately extend the Radon–Nikodym and Lebesgue decomposition theorems to the case where ν is a bounded signed measure (we keep the notation used in Section 7.3.2, so here µ remains positive!):

Theorem 7.20 Let µ be σ-finite (positive) measure and suppose that ν is a bounded signed measure. Then there is unique decomposition ν = νa + νs , into two signed measures, with νa µ and νs ⊥ µ. Moreover, R there is a unique (up to sets of µ-measure 0) h ∈ L1 (µ) such that νa (F ) = F h dµ for all F ∈ F.


213

Proof Given ν = ν + − ν − we wish to apply the Lebesgue decomposition and Radon– Nikodym theorems to the pairs of finite measures (ν + , µ) and (ν − , µ). First we need to check that for a signed measure λ µ we also have |λ| µ (for then clearly both λ+ µ and λ− µ). But if µ(E) = 0 and {Fi } partitions E, P then each µ(Fi ) = 0, hence λ(Fi ) = 0, so that i≥1 |λ(Fi )| = 0. As this holds for each partition, |λ(E)| = 0. Similarly, if λ is concentrated on a set A, and A ∩ E = Ø, then for any partition {Fi } of E we will have λ(Fi ) = 0 for every i ≥ 1. Thus |λ|(E) = 0, so |λ| is also concentrated on A. Hence if two signed measures are mutually singular, so are their total variation measures, and thus also their positive and negative variations. Applying the Lebesgue decomposition and Radon– Nikodym theorems to the measures ν + and ν − provides (positive) measures + + − − + + (ν that ν + = (ν + )a + R )0a , (ν )s , (ν− )a , (ν− )s such R (ν 00 )s , and (ν )a (F ) = − − h dµ, while ν = (ν )a + (ν )s and (ν )a (F ) = F h dµ, for non-negative F functions h0 , h00 ∈ L1 (µ), and with the measures (ν + )s , (ν − )s each mutually singular with µ. Letting νa = (ν + )a −(ν − )a we obtain R a signed measure νa µ, and a function h = h0 − h00 ∈ L1 (µ) with νa (F ) = F h dµ for all F ∈ F. The signed measure νs = (ν + )s − (ν − )s is clearly singular to µ, and h is unique up to µ-null sets, since this holds for h0 , h00 and the decomposition ν = ν + − ν − is minimal.

Example 7.5 R If g ∈ L1 (µ) then ν(E) = E g dµ is a signed measure and ν µ. The Radon– Nikodym theorem shows that (with our conventions) all signed measures ν µ have this form. We are nearly ready for the general form of the Fundamental Theorem of Calculus. First we confirm, as may be expected from the proof of the Radon– Nikodym theorem, the close relationship between the derivative of the bounded variation function F induced by a bounded signed (Borel) measure ν on R and the derivative f = F 0 :

Theorem 7.21 If ν is a bounded signed measure on R and F (x) = ν((−∞, x]) then for any a ∈ R, the following are equivalent:

(i) F is differentiable at a, and F 0 (a) = L.

ν(J) (ii) given ε > 0 there exists δ > 0 such that | m(J) − L| < ε if the open interval

214


J contains a and l(J) < ε.

Proof We may assume that L = 0; otherwise consider ρ = ν − Lm instead, restricted to a bounded interval containing a. If (i) holds with L = 0 and ε > 0 is given, we can find δ > 0 such that |F (y) − F (x)| < ε|y − x| whenever |y − x| < δ. Let J = (x, y) be an open interval containing a with (y −x) < δ. For sufficiently large N we can ensure that a > x + N1 > x and so for k ≥ 1, yk = x + N 1+k is bounded above by a and decreases to x as k → ∞. Thus |ν(yk , y])| = |F (y) − F (yk )| ≤ |F (y) − F (a)| + |F (a) − F (yk )| ≤ ε{(y − a) + (a − yk )} < εm(J).

ν(J) But since yk → x, ν(yk , y] → ν(x, y] and we have shown that | m(J) | < ε. Hence (ii) holds. For the converse, let ε, δ be as in (ii), so that with x < a < y and y − x < δ (ii) implies |ν(x, y + n1 )| < ε(y + n1 − x) for all large enough n. But T as (x, y] = n (x, y + n1 ), we also have

|ν(x, y]| < |F (y) − F (x)| < ε(y − x).

(7.3)

Finally, since (ii) holds, |ν({a})| ≤ |ν(I)| < εl(I) for any small enough open interval I containing a. Thus F (a) = F (a−) and so F is continuous at a. Since x < a < y < x + δ, we conclude that (7.3) holds with a instead of x, which shows that the right-hand derivative of F at a is 0, and with a instead of y which shows the same for the left-hand derivative. Thus F 0 (a) = 0, and so (i) holds.

Theorem 7.22 (Fundamental Theorem of Calculus) Let F be absolutely continuous on [a, b]. Then F is differentiable m-a.e. and its F Lebesgue–Stieltjes signed measure mF has Radon–Nikodym derivative dm dm = 0 F m-a.e. Moreover, for each x ∈ [a, b], Z x F (x) − F (a) = mF [a, x] = F 0 (t) dt. a


215

Proof dmF 1 RThe Radon–Nikodym theorem provides dm = h ∈ L (m) such that mF (E) = h dm for all E ∈ B. Choosing the partitions E

Pn = {(ti , ti+1 ] : ti = a +

i (b − a), i ≤ 2n } 2n

we obtain, successively, each Pn as the smallest common refinement of the partitions P1 , P2 , . . . , Pn−1 . Thus, setting hn (a) = 0 and n

hn (x) =

2 X mF (ti , ti+1 ] i=1

m(ti , ti+1 ]

n

1(ti ,ti+1 ] =

2 X F (ti+1 ) − F (ti ) i=1

ti+1 − ti

1(ti ,ti+1 ] for a < x ≤ b,

we obtain a sequence (hn ) corresponding to the sequence (hQn ) constructed in Step 2 of the proof of the Radon–Nikodym theorem. It follows that hn (x) → h(x) m-a.e. But for any fixed x ∈ (a, b), condition (ii) in Theorem 4.2 applied to the function F on each interval (ti , ti+1 ) with length less than δ, and with L = h(x), shows that h = F 0 m-a.e. The final claim is now obvious from the definitions. The following result is therefore immediate and it justifies the terminology ‘indefinite integral’ in this general setting.

Corollary 7.23 If F is absolutely continuous on [a, b] and F 0 = 0 m-a.e. then F is constant. A final corollary now completes the circle of ideas for distribution functions and their densities:

Corollary 7.24 Rx If f ∈ L1 ([a, b]) and F (x) = a f dm for each x ∈ [a, b] then F is differentiable m-a.e. and F 0 (x) = f (x) for almost every x ∈ [a, b]. As a further application of the Radon–Nikodym theorem we derive the Hahn–Jordan decomposition of ν which was outlined earlier. First we need the following

Theorem 7.25 Let ν be a bounded signed measure and let |ν| be its total variation. Then

216


we can find a measurable function h such that |h(ω)| = 1 for all ω ∈ Ω and R ν(E) = E h d|ν| for all E ∈ F.

Proof RThe Radon–Nikodym theorem provides a measurable function h with ν(E) = E h d|ν| for all E ∈ F since every |ν|-null set is ν-null ({E, Ø, Ø, . . .} is a partition of E). Let Cα = {ω : |h(ω)| < α} for α > 0. Then, for any partition {Ei } of Cα , X X X Z |ν(Ei )| = h d|ν| ≤ α|ν|(Ei ) = α|ν|(Cα ). i≥1

i≥1

Ei

i≥1

As this holds for any partition, it holds for their supremum, i.e. |ν|(Cα ) ≤ α|ν|(Cα ). For α < 1 we must conclude that Cα is |ν|-null, and hence also ν-null. Therefore |h| ≥ 1 ν-a.e. To show that |h| ≤ 1 ν-a.e. we note that if E has positive |ν|-measure, then, by definition of h, R h d|ν| |ν(E)| E = ≤ 1. |ν|(E) |ν|(E) That this implies |h| ≤ 1 ν-a.e. follows from the proposition below, applied with ρ = |ν|. Thus the set where |h| 6= 1 is |ν|-null, hence also ν-null, and we can redefine h there so that |h(ω)| = 1 for all ω ∈ Ω.

Proposition 7.26 Given a finite measure ρ and a function f ∈ L1 (ρ), suppose that for every R 1 E ∈ F with ρ(E) > 0 we have | ρ(E) E f dρ| ≤ 1. Then |f (ω)| ≤ 1, ρ-a.e. Hint Let E = {f > 1}. If ρ(E) > 0 consider

R

f E ρ(E)

dρ.

We are ready to derive the Hahn–Jordan decomposition very simply:

Proposition 7.27 Let ν be a bounded signed measure. There are disjoint measurable sets A, B such that A ∪ B = Ω and ν + (F ) = ν(B ∩ F ), ν − (F ) = ν(A ∩ F ) for all F ∈ F. Consequently, if ν = λ1 − λ2 for measures λ1 , λ2 then λ1 ≥ ν + and λ2 ≥ ν − . Hint Since dν = h d|ν| and |h| = 1 let A R = {h = −1}, B = {h = 1}. Use the definition of ν + to show that ν + (F ) = 21 F (1 + h) d|ν| = ν(F ∩ B) for every F.


217

Exercise 7.14 Let ν be a bounded signed measure. Show that for all F , ν + (F ) = supG⊂F ν(G), ν − (F ) = − inf G⊂F ν(G), all the sets concerned being members of F. Hint ν(G) ≤ ν + (G) ≤ ν(B ∩ G) + ν((B ∩ F )(B ∩ G)) = ν(B ∩ F ).

Exercise 7.15 R Show that when ν(F ) = F f dµ where f ∈ L1 (µ), where µ is a (positive) measure, the Hahn R decomposition sets are RA = {f < 0} and B = {f ≥ 0}, and ν + (F ) = F f + dν, while ν − (F ) = F f − dν. We finally arrive at a general definition of integrals relative to signed measures:

Definition 7.9 Let µ beR signed measure and f a measurable function on F ∈ F. Define the integral F f dµ by Z Z Z f dµ = f dµ+ − f dµ− F

F

F

whenever both terms on the right are finite or are not of the form ±(∞ − ∞). The function is sometimes called summable if the integral so defined is finite. Note that the earlier definition of a Lebesgue–Stieltjes signed measure fits into this general framework. We normally restrict attention to the case when both terms are finite, which clearly holds when µ is bounded.

Exercise 7.16 Verify the following: LetR µ be a finite measure and define the signed measure ν by ν(F ) = FRg dµ. Prove that f ∈ L1 (ν) if and only if R f g ∈ L1 (µ) and E f dν = E f g dµ for all µ-measurable sets E.

218


7.4 Probability 7.4.1 Conditional expectation relative to a σ-field Suppose we are given a random variable X ∈ L1 (P ), where (Ω, F, P ) is a probability space. In Chapter 5 we defined the conditional expectation E(X|G) of X ∈ L2 (P ) relative to a sub-σ-field G of F as the a.s. unique random variable Y ∈ L2 (G) satisfying the condition Z Z Y dP = X dP for all G ∈ G. (7.4) G

G

The construction was a consequence of orthogonal projections in the Hilbert space L2 with the extension to all integrable random variables undertaken ‘by hand’, which required a little care. With the Radon–Nikodym theorem at our disposal we can verify the existence of conditional expectations for integrable random variables very simply: R The (possibly signed) bounded measure ν(F ) = F X dP is absolutely continuous with respect to P . Restricting both measures to (Ω, G) maintains this relationship, so that Y R there is a G-measurable, P -a.s. unique random variable R such that ν(G) = G Y dP for every G ∈ G. But by definition ν(G) = G X dP , so the defining equation (7.4) of Y = E(X|G) has been verified.

Remark 7.2 In particular, this shows that for X ∈ L2 (F) its orthogonal projection onto 2 L R (G) is a version of the Radon–Nikodym derivative of the measure ν : F → X dP . F

We shall write E(X|G) instead of Y from now on, always keeping in mind that we have freedom to choose a particular ‘version’, i.e. as long as the results we seek demand only that relations concerning E(X|G) hold P -a.s., we can alter this random variable on a null set without affecting the truth of the defining equation:

Definition 7.10 A random variable E(X|G) is called the conditional expectation of X relative to a σ-field G if

(1) E(X|G) is G-measurable, R R (2) G E(X|G) dP = G X dP for all G ∈ G.


219

We investigate the properties of the conditional expectation. To begin with, the simplest are left for the reader as a proposition. In this and the subsequent theorem we make the following assumptions: (i) All random variables concerned are defined on a probability space (Ω, F, P );

(ii) X, Y and all (Xn ) used below are assumed to be in L1 (Ω, F, P );

(iii) G and H are sub-σ-fields of F.

The properties listed in the next proposition are basic, and are used time and again. Where appropriate we give verbal description of its ‘meaning’ in terms of information about X.

Proposition 7.28 The conditional expectation E(X|G) has the following properties: (i) E(E(X|G)) = E(X) (more precisely: any version of the conditional expectation of X has the same expectation as X). (ii) If X is G -measurable, then E(X|G) = X (if, given G, we already ‘know’ X, our ‘best estimate’ of it is perfect).

(iii) If X is independent of G, then E(X|G) = E(X) (if G ‘tells us nothing’ about X, our best guess of X is its average value).

(iv) (Linearity) E((aX + bY )|G) = aE(X|G) + bE(Y |G) for any real numbers a, b (note again that this is really says that each linear combination of versions of the right-hand-side is a version of the left-hand-side).

Theorem 7.29 The following properties hold for E(X|G) as defined above: (i) If X ≥ 0 then E(X|G) ≥ 0 a.s. (positivity). (ii) If {Xn }n≥1 are non-negative and increase a.s. to X, then {E(Xn |G)}n≥1 increase a.s. to E(X|G) (‘monotone convergence’ of conditional expectations). (iii) If Y is G-measurable and XY is integrable, then E(XY |G) = Y E(X|G) (‘taking out a known factor’). (iv) If H ⊂ G then E([E(X|G)]|H) = E(X|H) (the tower property).

220


(v) If ϕ : R → R is a convex function and ϕ(X) ∈ L1 (P ), then E(ϕ(X)|G) ≥ ϕ(E(X|G)). (This is known as the conditional Jensen inequality – a similar result holds for expectations. Recall that a real function ϕ is convex on (a, b) if for all x, y ∈ (a, b), ϕ(px + (1 − p)y) ≤ pϕ(x) + (1 − p)ϕ(y); the graph of ϕ stays on or below the straight line joining (x, ϕ(x)), (y, ϕ(y)).)

Proof (i) For each k ≥ 1 the set Ek = {E(X|G) < − k1 } ∈ G, so that Z Z X dP = E(X|G) dP. Ek

Ek

As X ≥ 0, the left-hand side is non-negative, while the right-hand-side is bounded above by − k1 P (Ek ). This forces P (Ek ) = 0 for each k, hence also S P (E(X|G) < 0} = P ( k Ek ) = 0. Thus E(X|G) ≥ 0 a.s. (ii) For each n let Yn be a version of E(Xn |G). By (i) and as in Section 5.4.3, the (Yn ) are non-negative and increase a.s. Letting Y = lim supn Yn provides a G-measurable random variable such that the real sequence (Y R n (ω))n converges to Y (ω) for almost all ω. Corollary 4.9 then shows that ( G Yn dP )n≥1 R R R increases to G Y dP . But we have G Yn dP = G Xn dP for each n, and (Xn ) increases pointwise to X. By Rthe monotone convergence theorem it follows that R R R ( G Xn dP )n≥1 increases to G X dP, so that G X dP = G Y dP. This shows that Y is a version of E(X|G) and therefore proves our claim. (iii) We can restrict attention to X ≥ 0, since the general case follows from this by linearity. Now first consider the case of indicators: if Y = 1E for some E ∈ G, we have, for all G ∈ G, Z Z Z Z 1E E(X|G) dP = E(X|G) dP = X dP = 1E X dP G

E∩G

E∩G

G

so that 1E E(X|G) satisfies the defining equation and hence is a version of the conditional expectation of the product XY. So E(XY |G) = Y E(X|G) has been verified when Y = 1E and E ∈ G. By the linearity property this extends to simple functions, and for arbitrary Y ≥ 0 we now use (ii) and a sequence (Yn ) of simple functions increasing to Y to deduce that, for non-negative X, E(XYn |G) = Yn E(X|G) increases to E(XY |G) on the one hand and to Y E(X|G) on the other. Thus if X and Y are both non-negative we have verified (iii). + − Linearity allows us R to extend this toR general Y = Y − Y . R R (iv) We have G E(X|G) dP = G X dP for G ∈ G and RH E(X|H) dP = H X dP for H ∈ H ⊂ G. Hence for H ∈ H we obtain H E(X|G) dP =


221

R

H E(X|H) dP. Thus E(X|H) satisfies the condition defining the conditional expectation of E(X|G) with respect to H, so that E[E(X|G)|H] = E(X|H). (v) A convex function can be written as the supremum of a sequence of affine functions, i.e. there are sequences (an ), (bn ) of reals such that ϕ(x) = supn (an x + bn ) for every x ∈ R. Fix n, then since ϕ(X(ω)) ≥ an X(ω) + bn for all ω, the positivity and linearity properties ensure that

E(ϕ(X)|G)(ω) ≥ E([an X + bn ]|G)(ω) = an E(X|G)(ω) + bn S for all ω ∈ ΩAn where P (An ) = 0. Since A = n An is also null, it follows that for all n ≥ 1, E(ϕ(X)|G)(ω) ≥ an E(X|G)(ω) + bn a.s. Hence the inequality also holds when we take the supremum on the right, so that (E(ϕ(X)|G)(ω) ≥ ϕ[(E(X|G)(ω)] a.s. This proves (v). An immediate consequence of (v) is that the LP -norm of E(X|G) is bounded by that of X for p ≥ 1, since the function ϕ(x) = |x|p is then convex: we obtain |E(X|G)|p = ϕ(E(X|G)) ≤ E(ϕ(X)|G) = E(|X|p |G) a.s. so that kE(X|G)kpp = E(|E(X|G)|p ) ≤ E(E(|X|p |G)) = E(|X|p ) = kXkpp , where the penultimate step applies property (1) to |X|p . Take pth roots to have kE(X|G)kp ≤ kXkp .

Exercise 7.17 Let Ω = [0, 1] with Lebesgue measure and let X(ω) = ω. Find E(X|G) if (a) G = {[0, 21 ], ( 12 , 1], [0, 1], Ø}, (b) G is generated by the family of sets {B ⊂ [0, 21 ], Borel}.

7.4.2 Martingales Suppose we wish to model the behaviour of some physical phenomenon by a sequence (Xn ) of random variables. The value Xn (ω) might be the outcome of the nth toss of a ‘fair’ coin which is tossed 1000 times, with ‘Heads’ recorded P as 1, ‘Tails’ as 0. Then Y (ω) = 1000 n=1 Xn (ω) would record the number of times that the coin had landed ‘Heads’. Typically, we would perform this random experiment a large number of times before venturing to make statements about the probability of ‘Heads’ for this coin. We could average our results, i.e. seek to compute E(Y ). But we might also be interested in guessing what the value of Xn (ω) might be after k < n tosses have been performed, i.e. for a fixed ω ∈ Ω,

222


does knowing the values of (Xi (ω))i≤k give us any help in predicting the value of Xn (ω) for n > k? In an ‘idealised’ coin-tossing experiment it is assumed that it does not, that is, the successive tosses are assumed to be independent — a fact which often perplexes the beginner in probability theory. There are many situations where the (Xn ) would represent outcomes where the past behaviour of the process being modelled can reasonably be taken to influence its future behaviour, e.g. if Xn records whether it rains on day n. We seek a mathematical description of the way in which our knowledge of past behaviour of (Xn ) can be codified. A natural idea is to use the σ-field Fk = σ{Xi : 0 ≤ i ≤ k} generated by the sequence (Xn )n≥0 as representing the knowledge gained from knowing the first k outcomes of our experiment. We call (Xn )n≥0 a (discrete) stochastic process to emphasise that our focus is now on the ‘dynamics’ of the sequence of outcomes as it unfolds. We include a 0th stage for notational convenience, so that there is a ‘starting point’ before the experiment begins, and then F0 represents our knowledge before any outcome is observed. So the information available to us by ‘time’ k (i.e. after k outcomes have been recorded) about the ‘state of the world’ ω is given by the values (Xi (ω))0≤i≤k and this is encapsulated in knowing which sets of Fk contain the point ω. But we can postulate a sequence of σ-fields (Fn )n≥0 quite generally, without reference to any sequence of random variables. Again, our knowledge of any particular ω is then represented at stage k ≥ 1 by knowing which sets in Fk contain ω. A simple example is provided by the binomial stock price model of Section 2.6.3 (see Exercise 2.13). Guided by this example, we turn this into a general

Definition 7.11 Given a probability space (Ω, F, P ) a (discrete) filtration is an increasing sequence of sub-σ-fields (Fn )n≥0 of F; i.e. F0 ⊂ F1 ⊂ F2 ⊂ . . . ⊂ Fn ⊂ . . . ⊂ F. We write F = (Fn )n≥0 . We say that the sequence (Xn )n≥0 of random variables is adapted to the filtration F if Xn is Fn -measurable for every n ≥ 0. The tuple (Ω, F, (Fn )n≥0 , P ) is called a filtered probability space. We shall normally assume in our applications that F0 = {Ø, Ω}, so that we begin with ‘no information’, and very often we shall assume that the ‘final’ σ-field generated by the whole sequence, i.e. F∞ = σ(∪n≥0 Fn ), is all of F (so that, by the end of the experiment, ‘we know all there is to know’). Clearly (Xn ) is adapted to its natural filtration (Fn )n , where Fn = σ(Xi : 0 ≤ i ≤ n}


223

for each n, and it is adapted to every filtration which contains this one. But equally, if Fn = σ(Xi : 0 ≤ i ≤ n} for some process (Xn )n , it may be that some for other process (Yn )n , each Yn is Fn -measurable, i.e. (Yn )n is adapted to (Fn )n . Recall that by Proposition 3.12 this implies that for each n ≥ 1 there is a Borel-measurable function fn : Rn+1 → R such that Yn = f (X0 , X1 , X2 , .., Xn ). We come to the main concept introduced in this section:

Definition 7.12 Let (Ω, F, (Fn )n≥0 , P ) be a filtered probability space. A sequence of random variables (Xn )n≥0 on (Ω, F, P ) is a martingale relative to the filtration F = (Fn )n≥0 provided: (i) (Xn )n is adapted to F; (ii) each Xn is in L1 (P ),

(iii) for each n ≥ 0, E(Xn+1 |Fn ) = Xn . We note two immediate consequences of this definition which are used over and over again: 1) If m > n ≥ 0 then E(Xm |Fn ) = Xn . This follows from the tower property of conditional expectations, since (a.s.) E(Xm |Fn ) = E(E(Xm |Fm−1 )|Fn ) = E(Xm−1 |Fn ) = . . . = E(Xn+1 |Fn ) = Xn . 2) Any martingale (Xn ) has constant expectation: E(Xn ) = E(E(Xn |F0 )) = E(X0 ) holds for every n ≥ 0, by 1) and (i) in Proposition 7.28. A martingale represents a ‘fair game’ in gambling: betting, for example, on the outcome of the coin tosses, our winnings in ‘game n’ (the outcome of the nth toss) would be ∆Xn = Xn − Xn−1 , that is the difference between what we had before and after that game. (We assume that X0 = 0.) If the games are fair we would predict at time (n − 1), before the nth outcome is known, that E(∆Xn |Fn−1 ) = 0, where Fk = σ{Xi : i ≤ k} are the σ-fields of the natural filtration of the process (Xn ). This follows because our knowledge at time (n − 1) is encapsulated in Fn−1 and in a fair game we would expect our incremental winnings at any stage to be 0 on average. Hence in this situation the (Xn ) form a martingale. Similarly, in a game favourable to the gambler we should expect that E(∆Xn |Fn−1 ) ≥ 0, i.e. E(Xn |Fn−1 ) ≥ Xn−1 a.s. We call a sequence satisfying this inequality (and (i), (ii) of Definition 7.12) a submartingale, while a

224


game unfavourable to the gambler (hence favourable to the casino!) is represented similarly by a supermartingale, which has E(Xn |Fn−1 ) ≤ Xn−1 a.s. for every n. Note that for a submartingale the expectations of the (Xn ) increase with n, while for a supermartingale they decrease. Finally, note that the properties of these processes do not change if we replace Xn by Xn − X0 (as long as X0 ∈ L1 (F0 ), to retain integrability and adaptedness) so that we can work without loss of generality with processes that start with X0 = 0.

Example 7.6 The most obvious, yet in some ways quite general, example of a martingale consists of a sequence of conditional expectations: given a random variable X ∈ L1 (F) and a filtration (Fn )n≥0 of sub-σ-fields of F, let Xn = E(X|Fn ) for every n. Then E(Xn+1 |Fn ) = E(E(X|Fn+1 )|Fn ) = E(X|Fn ) = Xn , using the tower property again. We can interpret this by regarding each Xn as giving us the information available at time n, i.e. contained in the σ-field Fn , about the random variable X. (Remember that the conditional expectation is the ‘best guess’ of X, with respect to mean-square errors, when we work in L2 .) For a finite filtration {Fn : 0 ≤ n ≤ N } with FN = F it is obvious that E(X|FN ) = X. For an infinite sequence we might hope similarly that ‘in the limit’ we will have ‘full’ information about X, which suggests that we should be able to retrieve X as the limit of the (Xn ) in some sense. The conditions under which limits exist require careful study — see e.g. [12], [8] for details. A second standard example of a martingale is:

Example 7.7 Suppose (Zn )n≥1 is a sequence of independent random variables with zero P mean. Let X0 = 0, F0 = {Ø, Ω}, set Xn = nk=1 Zk and define Fn = σ{Zk : k ≤ n} for each n ≥ 1. Then (Xn )n≥0 is a martingale relative to the filtration (Fn ). To see this recall that for each n, Zn is independent of Fn−1 , so that E(Zn |Fn−1 ) = E(Zn ) = 0. Hence E(Xn |Fn−1 ) = E(Xn−1 |Fn−1 ) + E(Zn ) = Xn−1 , since Xn−1 is Fn−1 -measurable. (You should check carefully which properties of the conditional expectation we used here!) A ‘multiplicative’ version of this example is the following:

Exercise 7.18 Let Zn ≥ 0 be a sequence of independent random variables with E(Zn ) =


225

µ = 1 Let Fn = σ{Zk : k ≤ n} and show that, X0 = 1, Xn = Z1 Z2 . . . Zn (n ≥ 1) defines a martingale for (Fn ), provided all the products are integrable random variables, which holds, e.g., if all Zn ∈ L∞ (Ω, F, P ).

Exercise 7.19 Let (Zn )n≥1 be a sequence of independent random variables with mean µ = E(Zn ) 6= 0 for all n. Show that the sequence of their partial sums Xn = Z1 + Z2 + · · · + Zn is not a martingale for the filtration (Fn )n , where Fn = σ{Zk : k ≤ n}. How can we ‘compensate’ for this by altering Xn ? Let X = (Xn )n≥0 be a martingale for the filtration F = (Fn )n≥0 (with our above conventions); briefly we simply refer to the martingale (X, F). The function φ(x) = x2 is convex, hence by Jensen’s inequality (Theorem 7.29) 2 we have E(Xn+1 |Fn ) ≥ (E(Xn+1 |Fn ))2 = Xn2 , so X 2 is a submartingale. We investigate whether it is possible to ‘compensate’, as in Exercise 7.19, to make the resulting process again a martingale. Note that the expectations of the X n2 are increasing, so we will need to subtract an increasing process from X 2 to achieve this. In fact, the construction of this ‘compensator’ process is quite general. Let Y = (Yn ) be any adapted process with each Yn ∈ L1 . For any process Z write its increments as ∆Zn = Zn − Zn−1 for all n. Recall that in this notation the martingale property can be expressed succinctly as E(∆Zn |Fn−1 ) = 0 — we shall use repeatedly in what follows. We define two new processes A = (An ) and M = (Mn ) with A0 = 0, M0 = 0, via their successive increments ∆An = E(∆Yn |Fn−1 ) and ∆Mn = ∆Yn − ∆An for n ≥ 1. We obtain E(∆Mn |Fn−1 ) = E([∆Yn − E(∆Yn |Fn−1 )]|Fn−1 ) = E(∆Yn |Fn−1 ) − E(∆Yn |Fn−1 ) = 0, as E(∆Yn |Fn−1 ) is Fn−1 -measurable. Hence M is a martingale. Moreover, the process A is increasing if and only if 0 ≤ ∆An = E(∆Yn |Fn−1 ) = E(Yn |Fn−1 ) − Yn−1 , which holds if and only if Y is a subPn Pn martingale. Note that An = k=1 ∆Ak = k=1 [E(Yk |Fk−1 ) − Yk−1 ] is Fn−1 measurable. Thus the value of An is ‘known’ by time n − 1. A process with this property is called predictable, since we can ‘predict’ its future values one step ahead. It is a fundamental property of martingales that they are not predictable: in fact, if X is a predictable martingale, then we have Xn−1 = E(Xn |Fn−1 ) = Xn a.s. for every n

226


where the first equality is the definition of martingale, while the second follows since Xn is Fn−1 -measurable. Hence a predictable martingale is a.s. constant, and if it starts at 0 it will stay there. This fact gives the decomposition of an adapted process Y into the sum of a martingale and a predictable process a useful uniqueness property: first, since M0 = 0 = A0 , we have Yn = Y0 + Mn + An for the processes M, A defined above. If also Yn = Y0 + Mn0 + A0n , where Mn0 is a martingale, and A0n is predictable, then Mn − Mn0 = A0n − An a.s. is a predictable martingale, 0 at time 0. Hence both sides are 0 for every n and so the decomposition is a.s. unique. We call this the Doob decomposition of an adapted process. It takes on special importance when applied to the submartingale Y = X 2 which arises from a martingale X. In that case, as we saw above, the predictable process A is increasing, so that An ≤ An+1 a.s. for every n, and the Doob decomposition reads: X 2 = X02 + M + A In particular, if X0 = 0 (as we can assume without loss of generality), we have written X 2 = M +A as the sum of a martingale M and a predictable increasing process A. The significance of this is revealed in a very useful property of martingales, which was a key component of the proof of the Radon–Nikodym theorem (see Step 1 (iv) of Theorem 7.3, where the martingale connection is well hidden!): for any martingale X we can write, with (∆Xn )2 = (Xn − Xn−1 )2 : 2 E(∆Xn )2 |Fn−1 ) = E([Xn2 − 2Xn Xn−1 + Xn−1 ]|Fn−1 )

2 = E(Xn2 |Fn−1 ) − 2Xn−1 E(∆Xn |Fn−1 ) − Xn−1 2 = E([Xn2 − Xn−1 ]|Fn−1 ).

Hence given the martingale X with X0 = 0, the decomposition X 2 = M + A yields, since M is also a martingale: 0 = E(∆Mn |Fn−1 ) = E((∆Xn )2 − ∆An )|Fn−1 ) 2 = E([Xn2 − Xn−1 ]|Fn−1 ) − E(∆An |Fn−1 ).

In other words, since A is predictable, E((∆Xn )2 |Fn−1 ) = E(∆An |Fn−1 ) = ∆An

(7.5)

which exhibits the process A as a conditional ‘quadratic variation’ process of the original martingale X. Taking expectations: E((∆Xn )2 ) = E(∆An ).


227

Example 7.8 Note also that E(Xn2 ) = E(Mn ) + E(An ) = E(An ) (why?), so that both sides are bounded for all n if and only if the martingale X is bounded as a sequence in L2 (Ω, F, P ). Since (An ) is increasing, the a.s. limit A∞ (ω) = limn→∞ An (ω) exists, and the boundedness of the integrals ensures in that case that E(A∞ ) < ∞.

Exercise 7.20 Suppose (Zn )n≥1 is a sequence of Bernoulli random variables, with each Zn taking the values 1 and −1, each with probability 12 . Let X0 = 0, Xn = Z1 +Z2 +· · ·+Zn , and let (Fn )n be the natural filtration generated by the (Zn ). Verify that (Xn2 ) is a submartingale, and find the increasing process (An ) in its Doob decomposition. What ‘unexpected’ property of (An ) can you detect in this example? In the discrete setting we now have the tools to construct ‘stochastic integrals’ and show that they preserve the martingale property. In fact, as we saw for Lebesgue–Stieltjes measures, for discrete distributions the ‘integral’ is simply an appropriate linear combination of increments of the distribution function. If we wish to use a martingale X as an integrator, we therefore need to deal with linear combinations of the increments ∆Xn = Xn − Xn−1 . Since we are now dealing with stochastic processes (that it, functions of both n and ω) rather than real functions, measurability conditions will help determine what constitutes an ‘appropriate’ linear combination. So, if for ω ∈ Ω we set I0 (ω) = 0 and form sums In (ω) =

n X k=1

ck (ω)(∆Xk )(ω) =

n X k=1

ck (ω)(Xk (ω) − Xk−1 (ω)) for n ≥ 1,

we look for measurability properties of the process (cn )n which ensure that the new process (In )n has useful properties. We investigate this when (cn )n is a bounded predictable process and X is a martingale for a given filtration (Fn )n . Some texts call the process (In )n a martingale transform — we prefer the term discrete stochastic integral. We calculate the conditional expectation of I n : E(In |Fn−1 ) = E([In−1 + cn ∆Xn ]|Fn−1 ) = In−1 + cn E(∆Xn |Fn−1 ) = In−1 , since cn is Fn−1 -measurable and E(∆Xn |Fn−1 ) = E(Xn |Fn−1 ) − Xn−1 = 0. Therefore, when the process c = (cn )n which is integrated against the martingale X = (Xn ), is predictable, the martingale property is preserved under the discrete stochastic integral: I = (In )n is also a martingale with respect to the

228


filtration (Fn )n . We shall write this stochastic integral as c · X, meaning that for all n ≥ 0, In = (c·X)n . The result has sufficient importance for us to record it as a theorem:

Theorem 7.30 Let (Ω, F, (Fn )n≥0 , P ) be a filtered probability space. If X is a martingale and c is a bounded predictable process, then the discrete stochastic integral c · X is again a martingale. Note that we use the boundedness assumption in order to ensure that ck ∆Xk is integrable, so that its conditional expectation makes sense. For L2 martingales (which are what we obtain in most applications) we can relax this condition and demand merely that cn ∈ L2 (Fn−1 ) for each n. While the preservation of the martingale property may please mathematicians, it is depressing news for gamblers! We can interpret the process c as representing the size of the stake the gambler ventures in every game, so that cn is the amount (s)he bets in game n. Note that cn could be 0, which mean that the gambler ‘sits out’ game n and places no bet. It also seems reasonable that the size of the stake depends on the outcomes of the previous games, hence cn is Fn−1 -measurable, and thus c is predictable. The conclusion that c · X is then a martingale means that ‘clever’ gambling strategies will be of no avail when the game is fair. It remains fair, whatever strategy the gambler employs! And, of course, if it starts out unfavourable to the gambler, so that X is a supermartingale (Xn−1 ≥ E(Xn |Fn−1 )), the above calculation shows that, as long as cn ≥ 0 for each n , then E(In |Fn−1 ) ≤ In−1 , so that the game remains unfavourable, whatever non-negative stakes the gambler places (and negative bets seem unlikely to be accepted, after all. . . ). You will verify immediately, of course, that a submartingale X produces a submartingale c · X when c is a non-negative process. Sadly, such favourable games are hard to find in practice. Combining the definition of (In )n with the Doob decomposition of the submartingale X 2 we obtain the identity which illustrates why martingales make useful ‘integrators’. We calculate the expected value of the square of (c · X)n when c = (cn ) is predictable and X = (Xn ) is a martingale: E((c ·

X)2n )

= E([

n X k=1

2

ck ∆Xk ] ) = E(

n X

cj ck ∆Xj ∆Xk ).

j,k=1

Consider terms in the double sum separately: when j < k we have E(cj ck ∆Xj ∆Xk ) = E(cj ck ∆Xj ∆Xk |Fk−1 ) = E(cj ck ∆Xj E(∆Xk |Fk−1 )) = 0


229

since the first three factors are all Fk−1 - measurable, while E(∆Xk |Fk−1 ) = 0 since X is a martingale. With j, k interchanged this also shows that these terms are 0 when k < j. The remaining terms have the form E(c2k (∆Xk )2 ) = E(c2k E((∆Xk )2 |Fk−1 )) = E(c2k ∆Ak ). By linearity, therefore, we have the fundamental identity for stochastic integrals relative to martingales (also called the Ito isometry): E([

n X k=1

ck ∆Xk ]2 ) = E(

n X

c2k ∆Ak ).

k=1

Remark 7.3 The sum inside the expectation sign on the right is a ‘Stieltjes sum’ for the increasing process, so that it is now at least plausible that this identity allows us to define martingale integrals in the continuous-time setting, using approximation of processes by simple processes, much as was done throughout this book for real functions. The Ito isometry is of critical importance in the definition of stochastic integrals relative to processes such as Brownian motion: in defining Lebesgue–Stieltjes integrals our integrators were of bounded variation. Typically, the paths of Brownian motion (a process we shall not define in this book — see (e.g.) [3] for its basic properties) are not of bounded variation, but the Ito isometry shows that their quadratic variation can be handled in the (much subtler) continuous-time version of the above framework, and this enables one to define integrals of a wide class of functions, using Brownian motion (and more general martingales) as ‘integrator’. We turn finally to the idea of stopping a martingale at a random time.

Definition 7.13 A random variable τ : Ω → {0, 1, 2, . . . , n, ..} ∪ {∞} is a stopping time relative to the filtration (Fn ) if for every n ≥ 1, the event {τ = n} belongs to Fn . Note that we include the value τ (ω) = ∞, so that we need {τ = ∞} ∈ F∞ = σ(∪n≥1 Fn ), the ‘limit σ-field’. Stopping times are also called random times, to emphasise that the ‘time’ τ is a random variable. For a stopping time τ the event {τ ≤ n} = ∪nk=0 {τ = k} is in Fn since for each k ≤ n {τ = k} ∈ Fk and the σ-fields increase with n. On the other hand,

230


given that for each n the event {τ ≤ n} ∈ Fn , then {τ = n} = {τ ≤ n}{τ ≤ n − 1} ∈ Fn . Thus we could equally well have taken the condition {τ ≤ n} ∈ Fn for all n as the definition of stopping time.

Example 7.9 A gambler may decide to stop playing after a random number of games, depending on whether his winnings X have reached a pre-determined level L (or his funds are exhausted!). The time τ = min{n : Xn ≥ L} is the first time at which the process X hits the interval [L, ∞); more precisely, for ω ∈ Ω, τ (ω) = n if Xn (ω) ≥ L while Xk (ω) < L for all k < n. Since {τ = n} is thus determined by the values of X and those of the Xk for k < n it is now clear that τ is a stopping time.

Example 7.10 Similarly, we may decide to sell our shares in a stock S if its value falls below 75% of its current (time 0) price. Thus we sell at the random time τ = min{n : Sn < 34 S0 }, which is again a stopping time. This is an example of a ‘stop-loss strategy’, and is much in evidence in a bear market. Quite generally, the first hitting time τA of a Borel set A ⊂ R by an adapted process X is defined by setting τA = min{n ≥ 0 : Xn ∈ A}. For any n ≥ 0 we have {τA ≤ n} = ∪k≤n {Xk ∈ A} ∈ Fn . To cater for the possibility that X never hits A we use the convention min Ø = ∞, so that {τA = ∞} = Ω(∪n≥0 {τA ≤ n}) represents this event. But its complement is in F∞ = σ(∪n≥0 Fn ), thus so is {τA = ∞}. We have proved that τA is a stopping time. Returning to the gambling theme, we see that stopping is simply a particular form of gambling strategy, and it should thus come as no surprise that the martingale property is preserved under stopping (with similar conclusions for super- and submartingales). For any adapted process X and stopping time τ , we define the stopped process X τ by setting Xnτ (ω) = Xn∧τ (ω) (ω) at each ω ∈ Ω. (Recall that for real x, y we write x ∧ y = min{x, y}.) The stopped process X τ is again adapted to the filtration (Fn )n , since {Xτ ∧n ∈ A} means that either τ > n and Xn ∈ A, or τ = k for some k ≤ n and Xk ∈ A. Now the event { Xn ∈ A} ∩ {τ > n} ∈ Fn , while for each k the event {τ = k} ∩ {Xk ∈ A} ∈ Fk . For all k ≤ n these events therefore all belong to Fn . Hence so does {Xτ ∧n ∈ A}, which proves that X τ is adapted.


231

Theorem 7.31 Let (Ω, F, (Fn ), P ) be a filtered probability space, and let X be a martingale with X0 = 0. If τ is a stopping time, the stopped process X τ is again a martingale.

Proof We use the preservation of the martingale property under discrete stochastic integrals (‘gambling strategies’). Let cn = 1{τ ≥n} for each n ≥ 1. This defines a bounded predictable process c = (cn ), since it takes only the values 0, 1 and {cn = 0} = {τ ≤ n − 1} ∈ Fn−1 , so that also {cn = 1} = Ω{cn = 0} ∈ Fn−1 . Hence by Theorem 7.30 the process c · X is again a martingale. But by construction (c · X)0 = 0 = X0 = X0τ , while for any n ≥ 1 (c · X)n = c1 (X1 − X0 ) + c2 (X2 − X1 ) + . . . + cn (Xn − Xn−1 ) = Xτ ∧n .

Since cn ≥ 0 as defined in the proof, it follows that the supermartingale and submartingale properties are also preserved under stopping. For a martingale we have, in addition, that expectation is preserved, i.e. (in general) E(Xτ ∧n ) = E(X0 ). Similarly, expectations increase for stopped submartingales, and decrease for stopped supermartingales. None of this, however, guarantees that the random variable Xτ defined by Xτ (ω) = Xτ (ω) (ω) has finite expectation – to obtain a result which relates its expectation to that of X0 we generally need to satisfy much more stringent conditions. For bounded stopping times (where there is a uniform upper bound N with τ (ω) ≤ N for all ω ∈ Ω), matters are simple: if X is a martingale, Xτ ∧n is integrable for all n, and by the above theorem E(Xτ ∧n ) = E(X0 ). Now apply this with n = N , so that Xτ ∧n = Xτ ∧N = Xτ . Thus we have E(Xτ ) = E(X0 ) whenever τ is a bounded stopping time. We shall not delve any further, but refer the reader instead to texts devoted largely to martingale theory, e.g. [12], [8], [3]. Bounded stopping times suffice for many practical applications, for example in the analysis of discrete American options in finance.

7.4.3 Applications to mathematical finance A major task and challenge for the theory of finance is to price assets and securities building models consistent with market practice. This consistency means that any deviation from the theoretical price should be penalized by the

232


market. Specifically, if a market player quotes a price different from the price provided by the model, she should be bound to lose money. The problem is the unknown future which has somehow to be reflected in the price since market participants express their beliefs about the future by agreeing on prices. Mathematically, an ideal situation is where the price process X(t) is a martingale. Then we would have the obvious pricing formula X(0) = E(X(T )) and in addition we would have the whole range of formulae for the intermediate prices by means of conditional expectation based on information gathered. However, the situation where the prices follow a martingale is incompatible with the market fact that money can be invested risk-free, which creates a benchmark for expectations for investors investing in risky securities. So we modify our task by insisting that the discounted values Y (t) = X(t) exp{−rt} form a martingale. The modification is by means of a deterministic constant so it does not create a problem for asset valuation. A particular goal is pricing derivative securities where we are given the terminal value (a claim) of the form f (S(T )), where f is known and the probability distributions of the values underlying asset S(t) are assumed to be known by taking some mathematical model. The above martingale idea would solve the pricing problem by constructing a process X(t) in such a way that X(T ) = f (S(T )) (We call X a replication of the claim.) We can summarise the tasks: Build a model of the prices of the underlying security S(t) such that 1. There is a replication X(t) such that X(T ) = f (S(T )), 2. The process Y (t) = X(t) exp{−rt} is a martingale, so Y (0) is the price of the security described by f, 3. Any deviation from the resulting prices leads to a loss. Steps 1 and 2 are mathematical in nature, but Step 3 is related to real market activities. We perform the task for the single step binomial model. Step 1. Recall that the prices in this model are S(0), S(1) = S(0)η where η = U or η = D with some probabilities. Let f (x) = (x − K)+ . We can easily find X(0) so that X(1) = (S(1) − K)+ by building a portfolio of n shares S and m units of bank account (see Section 6.5.5) after solving the system nS(0)U + mR = (S(0)U − K)+ ,

nS(0)D + mR = (S(0)D − K)+ . Note that X(1) is obtained by means of the data at time t = 0 and the model parameters. Step 2. Write R = exp{r}. The martingale condition we need is trivial: X(0)R = E(X(1)). The task here is to find the probability measure (X is


233

defined, we have no influence on R, so this is the only parameter we can adjust). Recall that we assume D ≤ R ≤ U . We solve X(0)R = pX(0)U + (1 − p)X(0)D which gives

R−D U −D and we are done. Hence the theoretical price of a call is p=

C = R−1 p(S(0)U − K) = exp{−r}E(S(1) − K)+ . Step 3. To see that within our model the price is right suppose that someone is willing to buy a call for C 0 > C. We sell it immediately investing the proceeds in the portfolio from step 1. At exercise our portfolio matches the call payoff and have earned the difference C 0 − C. So we keep buying the call until the seller realises the mistake and raises the price. Similarly if someone is selling a call at C 0 < C we generate cash by forming the (−n, −m) portfolio, buy a call (which, as a result of replication, settles our portfolio liability at maturity) and we have profit until the call prices quoted hit the theoretical price C. The above analysis summarises the key features of a general theory. A straightforward extension of this trivial model to n steps gives the so-called CRR (Cox–Ross–Rubinstein) model, which for large n is quite adequate for realistic pricing. We evaluated the expectation in the binomial model to establish the CRR price of a European call option in Proposition 4.37. For continuous time the task becomes quite difficult. In the Black–Scholes model S(t) = S(0) exp{(r − 12 σ 2 )t + σw(t)}, where w(t) is a stochastic process (called the Wiener process or Brownian motion) with w(0) = 0 and independent increments such that w(t)−w(s) is Gaussian with zero mean and variance t−s, s < t.

Exercise 7.21 Show that exp{−rt}S(t) is a martingale. Existence of the replication process X(t) is not easy but can be proved, as well as the fact that the process Y (t) is a martingale. This results in the same general pricing formula: exp{−rT }E(f (S(T )). Using the density of w(T ) this number can be written in an explicit form for particular derivative securities (see Section 4.7.5 where the formulae for the prices of a call and put options were derived).

234


Remark 7.4 The reader familiar with finance theory will notice that we are focused on pricing derivative securities and this results in considering the model where the discounted prices form a martingale. This model is a mathematical creation not necessarily consistent with real life, which requires a different probability space and different parameters within the same framework. The link between the real world and the martingale one is provided by a special application of Radon-Nikodym theorem which links the probability measures, but this is a story we shall not pursue here and refer the reader to numerous books devoted to the subject (for example [5]).

7.5 Proofs of propositions Proof (of Proposition 7.1) Suppose the (ε, δ)-condition fails. With A as in the hint, we have µ(A) ≤ S P∞ 1 1 µ(A) = 0. But µ( i≥n Fi ) ≤ i=n 2i = 2n−1 for every n ≥ 1. Thus S ν(Fn ) ≥ ε for every n, hence ν(En ) ≥ ε, where En = i≥n Fi . The sequence (En )n decreases, so that, as ν(F1 ) is finite, Theorem 2.13 (i) gives: ν(A) = ν(limn En ) = limn ν(En ) ≥ ε. Thus ν is not absolutely continuous with respect to µ. Conversely, if the (ε, δ)-condition holds, and µ(F ) = 0, then µ(F ) < δ for any δ > 0, and so, for every given ε > 0, ν(F ) < ε. Hence ν(F ) = 0. So ν µ.

Proof (of Proposition 7.4) We proceed as in Remark 4.1. Begin with R R g as the indicator function 1G for G ∈ G. Then we have: Ω g dµ = µ(G) = G hµ dϕ by construction of hµ . Next, Pn let g = i=1 ai 1Gi for sets G1 , G2 , . . . , Gn in G, and reals a1 , a2 , . . . , an ; then linearity of the integrals yields Z Z Z Z n n n X X X ghµ dϕ. g dµ = ai µ(Gi ) = ai ( hµ dϕ) = ai ( 1Gi hµ dϕ) = Ω

i=1

i=1

Gi

i=1

Ω

Ω

Finally, any G-measurable non-negative function g is approximated from below Pn by R an increasing sequence of G-simple functions gRn = i=1 ai 1Gi , its integral g dµ is the limit of the increasing sequence ( Ω gn hµ dϕ)n . But since 0 ≤ Ω hRµ ≤ 1 by construction, (gn hµ )Rincreases to ghµ pointwise, hence the sequence ( Ω gn hµ dϕ)n also increases to Ω ghµ dϕ, so the limits are equal. For integrable g = g + − g − apply the above to each of g + , g − separately, and use linearity.


235

Proof (of Proposition 7.6) S If n≥1 An = Ω and the An are not disjoint, replace them by En , where Sn−1 E1 = A1 , En = An ( i=1 Ei ), n > 1. The same can be done for the Bm and S hence we can take both sequences as disjoint. Now Ω = m,n≥1 (An ∩ Bm ) is also a disjoint union, and ν, µ are both finite on each An ∩ Bm . Re-order these sets into a single sequence (Cn )n≥1 and fix n ≥ 1. Restricting both measures to the σ-field Fn = {F ∩ Cn : F ∈ F} yields them as finite measures on (Ω, Fn ), so that the Radon-Nikodym theorem applies, R and provides a non-negative Fn measurable function hn such that ν(E) = E h dµ for each E ∈ Fn . But any S set F ∈ F has the form F = n Fn for Fn ∈ FnR, so we canRdefine h by P∞ setting h = hn for every n ≥ 1. Now ν(F ) = n=1 Fn hnR dµ = F h dµ. The uniqueness is clear: if g has the same properties as h,then F (h − g) dµ = 0 for each F ∈ F, so h − g = 0 a.e. by Theorem 4.15.

Proof (of Proposition 7.7) (i) This is trivial, since φ = λ + ν is σ-finite and absolutely continuous with respect to µ, and we have, for F ∈ F : Z Z dλ dν dφ dµ = φ(F ) = (λ + ν)(F ) = λ(F ) + ν(F ) = [ + ] dµ. dµ dµ dµ F F The integrands on the right and left extremes are thus a.s. (µ) equal, so the result follows. dν (ii) Write dλ dν = g and dµ = h. These are non-negative measurable functions and we need to show that, for F ∈ F Z λ(F ) = gh dµ. F

First consider this when g is replaced by a simple function of the form φ = Pn i=1 ci 1Ei . Then we obtain: Z

φ dν =

F

n X i=1

ci ν(F ∩ Ei ) =

n X i=1

ci

Z

h dµ =

F ∩Ei

Z

φh dµ.

F

Now let (φn ) be a sequence of simple functions increasing pointwise to g. Then by monotone convergence theorem: Z Z Z Z λ(F ) = g dν = lim φn dν = lim φn h dµ = gh dµ, F

n

F

n

F

since (φn h) increases to gh. This proves our claim.

F

236


Proof (of Proposition 7.8) Use the hint: λ1 + λ2 is concentrated on A1 ∪ A2 , while µ is concentrated on B1 ∩ B2 . But A1 ∪ A2 is disjoint from B1 ∩ B2 , hence the measures λ1 + λ2 and µ are mutually singular. This proves (i). For (ii), choose a set E for which µ(E) = 0 while λ2 is concentrated on E. Let F ⊂ E, so that µ(F ) = 0 and hence λ1 (F ) = 0 (since λ1 µ). This shows that λ1 is concentrated on E c , hence λ1 and λ2 are mutually singular. Finally, (ii) applied with λ1 = λ2 = ν shows that ν ⊥ ν, which can only happen when ν = 0.

Proof (of Proposition 7.11) Fix x ∈ R. The set A = {F (y) : y < x} is bounded above by F (x), while B = {F (y) : x < y} is bounded below by F (x). Hence sup A = K1 ≤ F (x), and inf B = K2 ≥ F (x) both exist in R and for any ε > 0, we can find y1 < x such that K1 − ε < F (y1 ) and y2 > x such that K2 + ε > F (y2 ). But since F is increasing this means that K1 − ε < F (y) < K1 throughout the interval (y1 , x) and K2 < F (y) < K2 + ε throughout (y, y2 ). Thus both one-sided limits F (x−) = limy↑x F (y) and F (x+) = lim y↓x F (y) are well-defined and by their definition F (x−) ≤ F (x) ≤ F (x+). Now let C = {x ∈ R : F is discontinuous at x}. For any x ∈ C we have F (x−) < F (x+). Hence we can find a rational r = r(x) in the open interval (F (x−), F (x+)). No two distinct x can have the same r(x), since if x1 < x2 we obtain F (x1 +) ≤ F (x2 −) from the definition of these limits. Thus the correspondence x ↔ r(x) defines a one-one correspondence between C and a subset of Q, so C is at most countable. At each discontinuity we have F (x−) < F (x+), so all discontinuities result from jumps of F.

Proof (of Proposition 7.12) Fix ε > 0, and let a finite set of disjoint intervals Jk = (xk , yk ) be given. Let S E = k Jk . Then n X

k=1

|F (yk ) − F (xk )| =

Z n X | k=1

yk

xk

f dm| ≤

n Z X k=1

yk

xk

|f | dm =

Z

E

|f | dm.

R But since f ∈ L1 , the measure µ(G) = G |f | dm is absolutely continuous with respect to Lebesgue measure m and hence, by Proposition 7.1, there exists δ > 0 such that m(F ) < δ implies µ(F ) < ε. But if the total length of the intervals Jk is less than δ, then m(F ) < δ, hence µ(F ) < ε. This proves that the function F is absolutely continuous.


237

Proof (of Proposition 7.14) Use the functions defined in the hint: for any partition (xk )k≤n of [a, x] we have F (x) − F (a) = =

n X

k=1 n X k=1

[F (xk ) − F (xk−1 )] [F (xk ) − F (xk−1 )]+ −

= p(x) − n(x)

n X

k=1

[F (xk ) − F (xk−1 )]−

so that p(x) = n(x)+[F (b)−F (a)] ≤ NF (x)+[F ((b)−F (a)] by definition of sup. This holds for all partitions, hence PF (x) = sup p(x) ≤ NF (x) + [F ((b) − F (a)]. On the other hand, writing n(x) = p(x)+[F (a)−F (b)] yields NF (x) ≤ PF (x)+ [F (a) − F (b)]. Thus PF (x) − NF (x) = F (b) − F (a). Now for any fixed partition we have TF (x) ≥

n X k=1

|F (xk ) − F (xk−1 )| = p(x) + n(x) = p(x) + {p(x) − [F (b) − F (a)]}

= 2p(x) − [PF (x) − NF (x)] = 2p(x) + NF (x) − PF (x). Take the supremum on the right: TF (x) ≥ 2PF (x) + NF (x) − PF (x) = Pn PF (x) + NF (x). But we can also write k=1 |F (xk ) − F (xk−1 )| = p(x) + n(x) ≤ PF (x) + NF (x) for any partition, so taking the sup on the left provides TF (x) ≤ PF (x) + NF (x). So the two sides are equal.

Proof (of Proposition 7.15) It will suffice to prove this for TF , as the other cases are similar. If the partition P of [a, b] produces the sum t(P) for the absolute differences, and if P 0 = P ∪{c} for some c ∈ (a, b), then t(P)[a, b] ≤ t(P 0 )[a, c] + t(P 0 )[c, b] and this is bounded above by TF [a, c] + TF [c, b] for all partitions. Thus it bounds TF [a, b] also. On the other hand, any partitions of [a, c] and [c, b] together make up a partition of [a, b], so that TF [a, b] bounds their joint sums. So the two sides must be equal. In particular, fixing a, TF ([a, c] ≤ TF [a, b] when c ≤ b, hence TF (x) = TF [a, x] is increasing with x. The same holds for PF and NF . The final statement is obvious.

Proof (of Proposition 7.17) Pn Pn (i) Given ε > 0, find δ > 0 such that i=1 (yi − xi ) < δ implies i=1 |F (yi ) − εδ F (xi )| < b−a . Given a partition (ti )i≤K of [a, b], we add further partition points, uniformly spaced and at a distance b−a M from each other, to ensure that

238


the combined partition (zi )i≤N has all its points less than δ apart. To do this we simply need to choose M as the integer part of T = b−a δ + 1. Since the (tj ) form a subset of the partition points (zi )i=0,1,...,N it follows that K X i=1

|F (ti ) − F (ti−1 )| ≤

N X i=1

|F (zi ) − F (zi−1 )|.

The latter sum can be re-ordered into M groups of terms where each group begins and ends with two consecutive new partition points: the k th group then contains (say) mk points altogether, and by their construction, the sum of their consecutive distances (i.e. the distance between the two new endpoints!) P mk εδ is less than δ, so for each k ≤ M, i=1 |F (wi,k ) − F (wi−1,k )| < b−a , where the (wi,k ) are the re-ordered points (zi ). Thus the whole sum is bounded by εδ εδ M ( b−a ) ≤ T ( b−a ) < ε. This shows that F ∈ BV [a, b], since the bound is independent of the original partition (ti )i≤K . For (ii), note first that, by (i), the function F has bounded variation on [a, b], so that over any subinterval [xi , yi ] the total variation function TF [xi , yi ] is finite. Again take ε, δ as given in the definition of absolutely continuPn ous functions. If (xi , yi )i≤n are subintervals with i=1 |yi − xi | < δ then Pn i=1 |F (yi ) − F (xi )| < ε. As in the previous proposition this implies that TF [xi , yi ] ≤ ε. Thus both PF [xi , yi ] and NF [xi , yi ] are less than ε, so that the functions F1 and F2 are absolutely continuous.

Proof (of Proposition 7.18) Obviously ν(Ø) = 0 and Ø ⊂ E for any E. So if ν is monotone increasing, ν(E) ≥ ν(Ø) ≥ 0. Hence ν is a measure. Conversely, if ν is a measure, F ⊂ E, ν(E) = ν(F ) + ν(E \ F ) ≥ ν(E).

Proof (of Proposition 7.26) R f f If E = {f > 1} has ρ(E) > 0, ρ(E) is well-defined and E ρ(E) dρ = R 1 f dρ > 1. This contradicts the hypothesis, so ρ(E) = 0. Similarly, ρ(E) E F = {f < −1} has ρ(F ) = 0. Hence |f | ≤ 1 ρ-a.e.

Proof (of Proposition 7.27) Choose h, A, B as in the hint. Recall that ν + = 12 (|ν| + ν), and note that 1 2 (1 + h) = h1B , so that, for F ∈ F, Z Z 1 + ν (F ) = (1 + h) d|ν|) = h d|ν| = ν(F ∩ B). 2 F F ∩B


239

But then, since B = Ac , ν − (F ) = −[ν(F ) − ν + (F )] = −[ν(F ) − ν(F ∩ B)] = −ν(F ∩ A). Finally, if ν = λ1 − λ2 where the λi are measures, then ν ≤ λ1 , so that ν + (F ) = ν(F ∩ B) ≤ λ1 (F ∩ B) ≤ λ1 (F ) by monotonicity. This proves the final statement of the proposition.

Proof (of Proposition 7.28) R R (i) Is immediate, as Ω E(X|G) dP = Ω X dP byRdefinition. R (ii) If both integrands are G-measurable and G E(X|G) dP = G X dP for all G ∈ G, then the integrands are a.s. equal by Theorem 4.15, and thus X is a version of E(X|G). (iii) For any G ∈ G, 1G and X are independent random variables, so that Z Z X dP = E(X1G ) = E(X)E(1G ) = E(X) dP G

G

Hence by definition E(X) is a version of E(X|G). But E(X) is constant, so the identity holds everywhere. (iv) Use the linearity of integrals: Z Z Z Z Z (aX + bY ) = a X dP + b Y dP = a E(X|G) dP + b E(Y |G) dP G G G G Z G = [aE(X|G) dP + bE(Y |G)] dP, G

so the result follows.

8 Limit theorems

In this chapter we introduce something of a change of pace and the reader may omit the more technically demanding proofs at a first reading, in order to gain an overview of the principal limit theorems for sequences of random variables. We put the spotlight firmly on probability to derive substantive applications of the preceding theory. First, however we discuss some basic modes of convergence of sequences of functions of real variable. Then we move to the probabilistic setting to which this chapter is largely devoted.

8.1 Modes of convergence Let E be a Borel subset of Rn . For a given sequence (fn ) in Lp (E), p ≥ 1, we can express the statement ‘fn → f as n → ∞’ in a number of distinct ways:

Definition 8.1 (1) fn → f uniformly on E: given ε > 0, there exists N = N (ε) such that, for all n ≥ N, kfn − f k∞ = sup (|fn (x) − f (x)|) < ε. x∈E

∞

(Note that we need fn ∈ L (E) for the sup to be finite in general.)

241

242


(2) fn → f pointwise on E: for each x ∈ E, given ε > 0, there exists N = N (ε, x) such that |fn (x) − f (x)| < ε for all n ≥ N . (3) fn → f almost everywhere (a.e.) on E: there is a null set F ⊂ E such that fn → f pointwise on E \ F .

(4) fn → f in Lp − norm (in the pth mean): kfn − f kp → 0 as n → ∞, i.e. for given ε > 0, ∃N = N (ε) such that kfn − fkp =

Z

p

E

|fn − f| dm

p1

0: Z 1 Z 1 1 1 |hn (x) − h(x)|p dx = xpn dx = xpn+1 → 0. pn + 1 0 0 0

Remark 8.1 There are still other modes of convergence which can be considered for sequences of measurable functions, and the relations between these and the above are quite complex in general. Here we will not pursue this theme in general, but specialize instead to probability spaces, where we derive additional relationships between the different limit processes.

Exercise 8.1 For each of the following decide whether fn → 0 (i) in Lp , (ii) uniformly, (iii) pointwise, (iv) a.e. (a) fn = 1[n,n+ n1 ] , (b) fn = n1[0, n1 ] − n1[− n1 ,0] .

8.2 Probability The remainder of this chapter is devoted to a discussion of the basic limit theorems for random variables in probability theory. The very definition of

244


‘probabilities’ relies on a belief in such results, i.e. that we can ascribe a meaning to the ‘limiting average’ of successes in a sequence of independent identically distributed trials. Then the purpose of the ‘endless repetition’ of tossing a coin is to use the ‘limiting frequency’ of Heads as the definition of the probability of Heads. Similarly, the pivotal role ascribed in statistics to the Gaussian density has its origin in the famous Central Limit Theorem (of which there is actually a large variety) which shows this density to describe the limit distribution of a sequence of distributions under appropriate conditions. Convergence of distributions therefore provides yet a further important limit concept for sequences of random variables. In both cases the concept of independence plays a crucial role and first of all we need to extend this concept to infinite sequences of random variables. In what follows, X1 , X 2 , . . . , X n , . . . will denote a sequence of random variables defined on a probability space.

Definition 8.2 We say that random variables X1 , X2 , . . . are independent if for any n ∈ N the variables X1 , . . . , Xn are independent (see Definition 3.3). An alternative is to demand that any finite collection of Xi be independent. Of course this condition implies independence since finite collections cover the initial segments of n variables. Conversely, take any finite collection of Xi and let n be the greatest index of this finite collection. Now X1 , . . . , Xn are independent and then for each subset the collection of its elements is independent; this includes, in particular, the chosen one. We study the following sequence Sn = X 1 + · · · + X n . If all Xi have the same distribution (we say that they are identically distributed ), then Snn is the average value of Xn (or X1 , it does not matter) after n repetitions of the same experiment. We study the bahaviour of Sn as n goes to infinity? The two main questions we address are: 1. When do the random variables Snn converge to a certain number? Here there is an immediate question of the appropriate mode of convergence. Positive answers to such questions are known as laws of large numbers.

8. Limit theorems

245

2. When do the distributions of the random variables Snn converge to a measure? Under what conditions is this limit measure Gaussian? The results we obtain in response are known as central limit theorems.

8.2.1 Convergence in probability Our first additional mode of convergence, convergence in probability, is sometimes termed ‘convergence in measure’.

Definition 8.3 A sequence X1 , X2 , . . . converges to X in probability if for each ε > 0 P (|Xn − X| > ε) → 0 as n → ∞.

Exercise 8.2 Go back to the proof of Theorem 8.1 (with E = [0, 1]) to see which of the sequences of random variables constructed there converge in probability.

Exercise 8.3 Find an example of a sequence of random variables on [0, 1] that does not converge to 0 in probability. We begin by showing that convergence almost surely (i.e. almost everywhere) is stronger than convergence in probability. But first we prove an auxiliary result.

Lemma 8.2 The following conditions are equivalent (a) Yn → 0 almost surely

(b) for each ε > 0,

lim P (

k→∞

∞ [

n=k

{ω : |Yn (ω)| ≥ ε}) = 0.

246


Proof Convergence almost surely, expressed succinctly, means that P ({ω : ∀ε > 0 ∃N ∈ N : ∀n ≥ N, |Yn (ω)| < ε}) = 1. Writing this set of full measure another way we have \ [ \ P( {ω : |Yn (ω)| < ε}) = 1. ε>0 N ∈N n≥N

The probability of the outer intersection (over all ε > 0) is less then the probability of any of its terms, but being already 1, it cannot increase, hence for all ε>0 [ \ P( {ω : |Yn (ω)| < ε}) = 1. N ∈N n≥N

We have a union of increasing sets so \ lim P ( {ω : |Yn (ω)| < ε}) = 1 N →∞

n≥N

thus lim (1 − P (

N →∞

\

n≥N

{ω : |Yn (ω)| < ε})) = 0

but we can write 1 = P (Ω) so that \ \ P (Ω) − P ( {ω : |Yn (ω)| < ε}) = P (Ω \ {ω : |Yn (ω)| < ε}) n≥N

n≥N

= P(

[

n≥N

{ω : |Yn (ω)| ≥ ε})

by De Morgan’s law. Hence (a) implies (b). Working backwards, these steps also prove the converse.

Theorem 8.3 If Xn → X almost surely then Xn → X in probability.

Proof For simplicity of notation consider the difference Yn = Xn −X and the problem reduces to the discussion of convergence of Yn to zero. We have lim P (

k→∞

∞ [

n=k

{ω : |Yn (ω)| ≥ ε}) ≥ lim P ({ω : |Yk (ω)| ≥ ε}) k→∞

and by Lemma 8.2 the limit on the left is zero hence so is that on the right.

8. Limit theorems

247

Note that the two sides of the inequality neatly summarize the difference between convergence a.s. and in probability. For convergence in probability we consider the probabilities that individual Yn are at least ε away from the limit, while for almost sure convergence we need to consider the whole tail sequence (Yn )n≥k simultaneously. The following example shows that the implication in the above theorem cannot be reversed and also shows that the convergence in Lp does not imply almost sure convergence.

Example 8.2 Consider the following sequence of random variables defined on Ω = [0, 1] with Lebesgue measure: Y1 = 1[0,1] , Y2 = 1[0,1/2] , Y3 = 1[1/2,1] , Y4 = 1[0,1/4] , Y5 = 1[1/4,1/2] and so on (like in the proof of Theorem 8.1). The sequence clearly converges to zero in probability and in Lp but for each ω ∈ [0, 1], Yn (ω) = 1 for infinitely many n, so it fails to converge pointwise. Convergence in probability has an additional useful feature:

Proposition 8.4 |X−Y | The function defined by d(X, Y ) = E( 1+|X−Y | ) is a metric and convergence in d is equivalent to convergence in probability.

R

Hint If Xn → X in probability then decompose the expectation into R + where A = {ω : |Xn (ω) − X(ω)| < ε}. A Ω\A

We now give a basic estimate of the probability of a non-negative random variable taking values in a given set by means of the moments of this random variable.

Theorem 8.5 (Chebyshev’s Inequality) If Y is a non-negative random variable, ε > 0, 0 < p < ∞, then P (Y ≥ ε) ≤

E(Y p ) . εp

(8.1)

Proof This is immediate from basic properties of integral: let A = {ω : Y (ω) ≥ ε}

248


and then E(Y p ) ≥

Z

Y p dP

(integration over a smaller set)

A

≥ P (A) · εp since Y p (ω) > εp on A, which gives the result after dividing by εp . Chebyshev’s inequality will be used mainly with small ε. But let us see what happens if ε is large.

Proposition 8.6 Assume that E(Y p ) < ∞. Then εp P (Y ≥ ε) → 0 as ε → ∞. Hint Write p

E(Y ) =

Z

p

Y dP + {ω:Y (ω)≥ε}

Z

Y p dP {ω:Y (ω) 0 and find δ > 0 such that if |x − y| < δ then |f (x) − f (y)| < (this is possible since f is uniformly continuous). Z Sn Sn Ex |f ( ) − f (x)| = |f ( ) − f (x)| dP n n {ω:| Snn −x| a) → 0 as a → ∞,

(8.2)

Sn − mn → 0 n

(8.3)

in probability.

We shall need the following lemma which is of interest in itself and will be useful in what follows.

Lemma 8.11 If Y ≥ 0, Y ∈ Lp , 0 < p < ∞, then Z ∞ E(Y p ) = py p−1 P (Y > y) dy. 0

In particular (p = 1) E(Y ) =

Z

∞ 0

P (Y > y) dy.

8. Limit theorems

253

Proof (of the Lemma) This is a simple application of Fubini’s theorem: Z ∞ py p−1 P (Y > y) dy 0 Z ∞Z = py p−1 1{ω:Y (ω)>y} (ω) dP (ω) dy Ω Z0 Z ∞ = py p−1 1{ω:Y (ω)>y} (ω) dy dP (ω) (by Fubini) Ω

0

Z Z

=

ZΩ

=

Y (ω)

py p−1 dy dP (ω)

0

Y p (ω) dP (ω) (computing the inner integral, ω fixed)

Ω

= E(Y p ) as required.

Proof (of the Theorem) Take ε > 0 and obviously P (|

Sˆn Sn − mn | ≥ ε) ≤ P (| − mn | ≥ ε) + P (Sˆn 6= Sn ). n n

We estimate the first term P (|

ˆ E(| Snn − mn |2 ) Sˆn − mn | ≥ ε) ≤ (by Chebyshev) 2 n Pε E(| Xk (n) − nmn |2 ) = 2 2 P n ε Var( Xk (n)) = . n2 ε 2

Note that the truncated random variables are independent, being functions of

254


the original ones, hence we may continue the estimation: P Var(Xk (n)) ≤ k n2 ε 2 Var(X1 (n)) = (as Var(Xk (n)) are the same) nε2 E(X12 (n)) (by Var(Z) = E(Z 2 ) − (EZ)2 ≤ E(Z 2 )) ≤ 2 nε Z ∞ 1 2yP (|X1 (n)| > y) dy (by the Lemma for p = 2) = 2 nε 0 Z n 1 = 2 2yP (|X1 | > y) dy. nε 0 The function y 7→ 2yP (|X1 | > y) converges to 0 as y → ∞ by hypothesis, hence for given δ > 0 there is y0 such that for y ≥ y0 this quantity is less than 1 2 2 δε , and we have Z y0 Z n 1 1 = 2 2yP (|X1 | > y) dy + 2 2yP (|X1 | > y) dy nε 0 nε y0 1 1 δε2 y max {yP (|X | > y)} + n 0 1 nε2 y∈[0,∞] nε2 2 ≤δ

≤

provided n is sufficiently large. So the first term converges to 0. Now we tackle the second term: P (Sˆn 6= Sn ) ≤ P (Xk (n) 6= Xk for some k ≤ n) n X ≤ P (Xk (n) 6= Xk ) (by subadditivity of P ) k=1

= nP (X1 (n) 6= X1 )

(the same distributions)

= nP (|X1 | > n)

→ 0 by hypothesis. This completes the proof.

Remark 8.3 Note that we cannot generalize the last theorem to the case of uncorrelated random variables since we made essential use of the independence. Although the identity X X Var Xk (n) = Var(Xk (n))

8. Limit theorems

255

holds for uncorrelated random variable, we needed the independence of the (Xk ) — which implies that of the (Xk (n)) — to conclude that the truncated random variables are uncorrelated.

Theorem 8.12 If Xn are independent and identically distributed, E(|X1 |) < ∞, then (8.2) is satisfied, mn → m = E(X1 ) and Snn → m in probability. (Note that we do not assume here that X1 has finite variance.)

Proof The finite expectation of |X1 | gives condition (8.2): Z aP (|X1 | > a) = a 1{ω:|X1 (ω)|>a} dP Z Ω ≤ |X1 |1{ω:|X1 (ω)|>a} dP ZΩ = |X1 | dP →0

{ω:|X1 (ω)|>a}

as a → ∞ by the dominated convergence theorem. Hence Snn − mn → 0 but mn = E(X1 1{ω:|X1 (ω)|≤n} ) → E(X1 ) as n → ∞ so the result follows.

8.2.3 The Borel–Cantelli Lemmas The idea that a point ω ∈ Ω belongs to ‘infinitely many’ events of a given sequence (An ) ⊂ F can easily be made precise: for every n ≥ 1 we need to be able to find an mn ≥ n such that ω ∈ Amn . This identifies a subsequence S (mn ) of indices such that for each n ≥ 1, ω ∈ Amn , i.e. ω ∈ m≥n Am for every n ≥ 1. Thus we say that ω ∈ An infinitely often, and write ω ∈ An i.o., T∞ S∞ if ω ∈ n=1 m=n Am . We call this set the upper limit of the sequence (An ) and write it as ∞ [ ∞ \ lim sup An = Am . n→∞

n=1 m=n

Exercise 8.6 Find lim supn→∞ An for a sequence A1 = [0, 1], A2 = [0, 21 ], A3 = [ 21 , 1], A4 = [0, 41 ], A5 = [ 41 , 21 ] etc.

256


Given ε > 0, a sequence of random variables (Xn ) and a random variable X, for each n set An = {|Xn − X| > ε}. Then ω ∈ An i.o. precisely when for every ε > 0, |Xn (ω) − X(ω)| > ε occurs for all elements of an infinite subsequence (mn ) of indices, which means that (Xn ) fails to converge to X. Hence it follows that Xn −→ X a.s.(P )

⇐⇒

∀ε > 0 P (lim sup An ) = P (|Xn −X| > ε i.o.) = 0. n→∞

Similarly, define lim inf An = n→∞

∞ \ ∞ [

An .

n=1 m=n

(We say that this set is the lower limit of the sequence (An ).)

Proposition 8.13 (i) We have ω ∈ lim inf n→∞ An if and only if ω ∈ An except for finitely many n. (We say that ω ∈ An eventually.)

(ii) P (Xn −→ X) = limε→0 P (|Xn − X| < ε eventually).

(iii) If A = lim supn→∞ An then Ac = lim inf n→∞ Acn .

(iv) For any sequence (An ) of events, P ({ω ∈ An ev.}) ≤ P ({ω ∈ An i.o.}) Hint: Use Fatou’s lemma on the indicator functions of the sets in (iv). The sets lim inf n→∞ An and lim supn→∞ An are ‘tail events’ of the sequence (An ) : we can only determine whether a point belongs to them by knowing the whole sequence. It is frequently true that the probability of a tail event is either 0 or 1 - such results are known as 0 − 1 laws. The simplest of these is provided by combining the two ‘Borel-Cantelli lemmas’ to which we now turn: together they show that for a sequence of independent events (An ), lim supn→∞ An has either probability 0 or 1, depending on whether the series of their individual probabilities converges or diverges. In the first case, we do not even need independence, but can prove the result in general.

Exercise 8.7 Let Sn = X1 + X2 + ... + Xn describe the position after n steps of a symmetric random walk on Zd . Using the asymptotic formula: n! ∼ √ n n 2πn and the Borel-Cantelli lemmas show that the probability of e {Sn = 0 i.o.} is 1 when d = 1, 2 and 0 for d > 2. We have the following simple but fundamental fact.

8. Limit theorems

257

Theorem 8.14 (Borel–Cantelli Lemma) If ∞ X

n=1

then

P (An ) < ∞

P (lim sup An ) = 0 n→∞

i.e. ‘ω ∈ An for infinitely many n’ occurs only with probability zero.

Proof First note that lim supn→∞ An ⊂

S∞

P (lim sup An ) ≤ P ( n→∞

≤

n=k

An hence

∞ [

An ) (for all k)

n=k

∞ X

P (An ) (by subadditivity)

n=k

→0

since the tail of a convergent series converges to 0. The basic application of the lemma provides a link between almost sure convergence and convergence in probability.

Theorem 8.15 If Xn → X in probability then there is a subsequence Xkn converging to X almost surely.

Proof We have to find a set of full measure on which a subsequence would converge. So the set on which the behaviour of the whole sequence is ‘bad’ should be of measure zero. For this we employ the Borel–Cantelli lemma whose conclusion is precisely that. So we introduce the sequence An encapsulating the ‘bad’ behaviour of Xn , which from the point of convergence is expressed by inequalities of the type |Xn (ω) − X(ω)| > a. Specifically, we take a = 1 and since P (|Xn − X| > 1) → 0 we find k1 such that P (|Xn − X| > 1) < 1

258


for n ≥ k1 . Next for a =

1 2

we find k2 > k1 such that for all n ≥ k2 P (|Xn − X| >

1 1 )< . 2 4

We continue that process obtaining an increasing sequence of integers kn with P (|Xkn − X| >

1 1 ) < 2. n n

P∞ The series n=1 P (An ) converges, where An = {ω : |Xkn (ω) − X(ω)| > n1 }, hence A = lim sup An has probability zero. We observe that for ω ∈ Ω \ A, lim sup Xkn (ω) = X(ω). For, if ω ∈ Ω \ A, T 1 then for some k, ω ∈ ∞ n=k (Ω \ An ) so for all n ≥ k, |Xkn (ω) − X(ω)| ≤ n , hence we have obtained the desired convergence. The second Borel–Cantelli lemma partially completes the picture. Under the additional condition of independence it shows when the probability that infinitely many events occur is one.

Theorem 8.16 Suppose that the events An are independent. We have ∞ X

n=1

P (An ) = ∞

⇒

P (lim sup An ) = 1. n→∞

Proof It is sufficient to show that for all k P(

∞ [

An ) = 1

n=k

since then the intersection over k will also have probability 1. Fix k and consider the partial union up to m > k. The complements of An are also independent hence m m m \ Y Y P( Acn ) = P (Acn ) = (1 − P (An )). n=k

Since 1 − x ≤ e

−x

m Y

n=k

n=k

n=k

,

(1 − P (An )) ≤

m Y

n=k

e

−P (An )

= exp(−

m X

n=k

P (An )).

8. Limit theorems

259

The last expression converges to 0 as m → ∞ by the hypothesis, hence P(

m \

n=k

but P(

m \

Acn ) = P (Ω \

n=k Sm n=k An

Acn ) → 0

m [

n=k

An ) = 1 − P (

m [

An ).

n=k

S S∞ The sets Bm = form an increasing chain with ∞ A m=k Bm = S∞n=k n and so P (Bm ), which as we know converges to 1, converges to P ( n=k An ). Thus this quantity is also equal to 1. Below we discuss strong laws of large numbers, where convergence in probability is strengthened to almost sure convergence. But already we can observe some limitations of these improvements. Drawing on the second Borel–Cantelli lemma we give a negative result.

Theorem 8.17 Suppose that X1 , X2 , . . . are independent identically distributed random variables and assume that E(|X1 |) = ∞ (hence also E(|Xn |) = ∞ for all n). Then

(i) P ({ω : |Xn (ω)| ≥ n for infinitely many n}) = 1,

(ii) P (limn→∞

Sn n

exists and is finite) = 0.

Proof (i) First E(|X1 |) = = ≤

Z

∞

P (|X1 | > x) dx

0 ∞ Z k+1 X

k=0 ∞ X

k=0

k

(by Lemma 8.11)

P (|X1 | > x) dx

(countable additivity)

P (|X1 | > k)

because the function x 7→ P (|X1 | > x) reaches its maximum on [k, k + 1] for x = k since {ω : |X1 (ω)| > k} ⊃ {ω : |X1 (ω)| > x} if x ≥ k. By the hypothesis this series is divergent, but P (|X1 | > k) = P (|Xk | > k) as the distributions are identical, so ∞ X P (|Xk | > k) = ∞. k=0

260


The second Borel–Cantelli lemma is applicable yielding the claim. (ii) Denote by A the set where the limit of elementary algebra of fractions gives

Sn n

exists (and is finite). Some

Sn Sn+1 (n + 1)Sn − nSn+1 Sn − nXn+1 Sn Xn+1 − = = = − . n n+1 n(n + 1) n(n + 1) n(n + 1) n+1 For any ω0 ∈ A the left-hand side converges to zero and also Sn (ω0 ) → 0. n(n + 1) Hence also

Xn+1 (ω0 ) n+1

→ 0. This means that

ω0 ∈ / {ω : |Xk (ω)| > k for infinitely many k} = B, say, so A ⊂ Ω \ B. But P (B) = 1 by (i), hence P (A) = 0.

8.2.4 Strong law of large numbers We shall consider several versions of the strong law of large numbers, first by imposing additional conditions on the moments of the sequence (Xn ), and then gradually relaxing these we arrive at Theorem 8.21, which provides the most general positive result. The first result is due to von Neumann. Note that we do not impose the condition that the Xn have identical distributions. The price we pay is having to assume that higher order moment are finite. However, for many familiar random variables, Gaussian for example, this is not a serious restriction.

Theorem 8.18 Suppose that the random variables Xn are independent, E(Xn ) = m, and E(Xn4 ) ≤ K. Then n Sn 1X = Xk → m a.s. n n k=1

Proof By considering Xn −m we may assume that E(Xn ) = 0 for all n. This simplifies

8. Limit theorems

261

the following computation E(Sn4 )

n X

=E

Xk

k=1



= E

n X

!4

Xk4 +

k=1

X

+

i,j,k,l distinct

X

Xi2 Xj2 +

i6=j

X

Xi Xj Xk2

i6=j



Xi Xj Xk Xl  .

The last two terms vanish by independence: X X X E( Xi Xj Xk2 ) = E(Xi Xj Xk2 ) = E(Xi )E(Xj )E(Xk2 ) = 0 i6=j

i6=j

i6=j

and similarly for the term with all indices distinct X X E( Xi Xj Xk Xl ) = E(Xi Xj Xk Xl ) X = E(Xi )E(Xj )E(Xk )E(Xl ) = 0.

The first term is easily estimated by the hypothesis X X E( Xk4 ) = E(Xk4 ) ≤ nK.

To the remaining term we first apply the Schwarz inequality q X X Xq E(Xi4 ) E(Xj4 ) ≤ N K E( Xi2 Xj2 ) = E(Xi2 Xj2 ) ≤ i6=j

i6=j

i6=j

where N is the number of components of this kind. (We could do better by employing independence, but then we would have to estimate the second moments by the fourth and it would boil down to the same.) N first note that the pairs of two distinct indices can be chosen in To find n(n−1) n ways. Having fixed, i, j the term Xi2 Xj2 arises in 6 ways corre= 2 2 sponding to possible arrangements of 2 pairs of 2 indices: (i, i, j, j), (i, j, i, j), (i, j, j, i), (j, j, i, i), (j, i, j, i), (j, i, i, j). So N = 3n(n − 1) and we have E(Sn4 ) ≤ K(n + 3n(n − 1)) = K(n + 3n2 − 3n) ≤ 3Kn2 . By Chebyshev’s inequality P (|

Sn E(Sn4 ) 3K 1 | ≥ ε) = P (|Sn | ≥ nε) ≤ ≤ 4 · 2. 4 n (nε) ε n

P 1 The series n2 converges and by Borel–Cantelli the set lim sup An with An = {ω : | Snn | ≥ ε} has measure zero. Its complement is the set of full measure we

262


need on which the sequence Snn converges to 0. To see this let ω ∈ / lim sup An which means that ω is in finitely many An . So for a certain n0 , all n ≥ n0 , ω∈ / An , i.e. Snn < ε (as observed before). and this is precisely what was needed for the convergence in question. The next law will only require finite moments of order 2, even not necessarily uniformly bounded. We precede it by an auxiliary but crucial inequality due to Kolmogorov. It gives a better estimate than does the Chebyshev inequality. The latter says that Var(Sn ) . P (|Sn | ≥ ε) ≤ ε2 In the theorem below the left-hand side is larger hence the result is stronger.

Theorem 8.19 If X1 , . . . , Xn are independent with 0 expectation and finite variances, then for any ε > 0 Var(Sn ) P ( max |Sk | ≥ ε) ≤ 1≤k≤n ε2 where Sn = X1 + · · · + Xn .

Proof We fix an ε > 0 and describe the first instance that |Sk | exceeds ε. Namely, we write 1 if |S1 | < ε, . . . , |Sk−1 | < ε, |Sk | ≥ ε ϕk = 0 if all |Si | < ε. For any ω at most one of the numbers ϕk (ω) may be 1, the remaining ones being 0, hence their sum is either 0 or 1. Clearly n X

ϕk = 0

⇔

1≤k≤n

n X

ϕk = 1

⇔

1≤k≤n

k=1

k=1

Hence

P ( max |Sk | ≥ ε) = P ( 1≤k≤n

max |Sk | < ε, max |Sk | ≥ ε.

n X k=1

ϕk = 1) = E(

n X k=1

ϕk )

8. Limit theorems

263

since the expectation is the integral of an indicator function: Z Z n X E( ϕk ) = 0 dP + P P {ω:

k=1

n k=1

ϕk (ω)=0}

{ω:

n k=1

1 dP.

ϕk (ω)=1}

So it remains to show that E(

n X k=1

ϕk ) ≤

1 1 Var(Sn ) = 2 E(Sn2 ), 2 ε ε

the last equality because E(Sn ) = 0. We estimate E(Sn2 ) from below E(Sn2 ) ≥ E( = E( ≥ E( = E(

n X

k=1 n X

k=1 n X k=1 n X

n X

ϕk · Sn2 ) (since

k=1

ϕk ≤ 1)

[Sk2 + 2Sk (Sn − Sk ) + (Sn − Sk )2 ]ϕk ) (simple algebra) [Sk2 + 2Sk (Sn − Sk )]ϕk ) Sk2 ϕk ) + 2E(

k=1

n X k=1

(non-negative term deleted)

Sk (Sn − Sk )ϕk ).

We show that the last expectation is equal to 0. Observe that ϕk is a function of random variables X1 , . . . , Xk so it is independent of Xk+1 , . . . , Xn and also Sk is independent of Xk+1 , . . . , Xn for the same reason. We compute one component of this last sum: E(Sk (Sn − Sk )ϕk ) = E(Sk ϕk (

n X

Xi ))

i=k+1 n X

= E(Sk ϕk )E(

(by the definition of Sn )

Xi ) (by independence)

i=k+1

= 0 (since E(Xi ) = 0 for all i). In the remaining sum note that for each k ≤ n, ϕk Sk2 ≥ ϕk ε2 (this is 0 ≥ 0 if ϕk = 0 and Skk ≥ ε2 if ϕk = 1, both true), hence E(

n X k=1

Sk2 ϕk ) ≥ E(ε2

which gives the desired inequality.

n X k=1

ϕk ) = ε2 E(

n X k=1

ϕk )

264


Theorem 8.20 Suppose that X1 , X2 , . . . are independent with E(Xn ) = 0 and ∞ X 1 Var(Xn ) < ∞. n2 n=1

Then

Sn →0 n

almost surely.

Proof We introduce auxiliary random variables Ym = max |Sk | m k≤2

and for 2m−1 ≤ n ≤ 2m |

Sn 1 1 | ≤ max |Sk | ≤ m−1 Ym . n n k≤2m 2

m It is sufficient to show that Y2m → 0 almost surely and by Lemma 8.2 it is sufficient to show that for each ε > 0

∞ X

m=1

P (|

Ym | ≥ ε) < ∞. 2m

First take a single term P (|Ym | ≥ 2m ε) and estimate it by Kolmogorov’s inequality (Theorem 8.19) P (|Ym | ≥ 2m ε) ≤

Var(S2m ) . ε2 22m

The problem reduces to showing that ∞ X

m=1

Var(S2m )

1 < ∞. 4m

8. Limit theorems

265

We rearrange the components m

∞ 2 X 1 1 X Var(Xk ) Var(S2m ) m = 4 4m m=1 m=1 ∞ X

k=1 ∞ X

= Var(X1 )

∞ X 1 1 + Var(X ) 2 m m 4 4 m=1 m=1

+Var(X3 ) +Var(X5 ) +···

∞ ∞ X X 1 1 + Var(X ) 4 m m 4 4 m=2 m=2

∞ ∞ X X 1 1 + · · · + Var(X ) 8 m m 4 4 m=3 m=3

since Var(X1 ), Var(X2 ) appear in all components of the series in m (1, 2 ≤ 21 ), Var(X3 ), Var(X4 ) appear in all except the first one (21 < 3, 4 ≤ 22 ), Var(X5 ), . . . , Var(X8 ) appear in all except the first two (22 < 5, 6, 7, 8 ≤ 23 ), and so on. We arrive at the series ∞ X

Var(Xk )ak

k=1

where ak =

X

{m:2m >k}

1 . 4m

This is a geometric series with ratio 14 and the first term 41j where j is the least integer such that 2j > k. If we replace 2j by k we increase the sum by adding more terms and the first terms is then 21k : ak ≤

1 2k

1−

1 4

and by the hypothesis ∞ X

k=1

Var(Xk )ak ≤

∞

4X 1 Var(Xk ) k < ∞ 3 2 k=1

which completes the proof. Finally, we relax the conditions on moments even further; simultaneously we need to impose the assumption that the random variables are identically distributed.

266


Theorem 8.21 Suppose that X1 , X2 , . . . are independent identically distributed, with E(X1 ) = m < ∞. Then Sn → m almost surely. n

Proof The idea is to use the previous theorem where we needed finite variances. Since we do not have that here we truncate Xn Yn = Xn (n) = Xn 1{ω:|Xn (ω)|≤n} . The truncated random variables have finite variances since each is bounded: |Yn | ≤ n (the right-hand side is forced to be zero if Xn dare upcross the level n). The new variables differ from the original ones if |Xn | > n. This, however, cannot happen too often as the following argument shows. First ∞ X

n=1

= = ≤ ≤

P (Yn 6= Xn )

∞ X

n=1 ∞ X

P (|Xn | > n) P (|X1 | > n) (the distributions being the same)

n=1 ∞ Z n X

n=1 n−1 Z ∞ 0

P (|X1 | > x) dx

(as P (|X1 | > x) ≥ P (|X1 | > n))

P (|X1 | > x) dx

= E(|X1 |) (by Lemma 8.11) < ∞.

So by Borel–Cantelli, with probability one only finitely many events Xn 6= Yn happen, so in other words, there is a set Ω 0 with P (Ω 0 ) = 1 such that for n ω ∈ Ω 0 , Xn (ω) = Yn (ω) for all except finitely many n. So on Ω 0 if Y1 +···+Y n Sn converges to some limit, then the same holds for n . To use the previous theorem we have to show the convergence of the series ∞ X Var(Yn ) n2 n=1

8. Limit theorems

267

but since Var(Yn ) = E(Yn2 ) − (EYn )2 ≤ EYn2 it is sufficient to show the convergence of ∞ X E(Yn2 ) . n2 n=1 To this end note first Z ∞ E(Yn2 ) = 2xP (|Yn | > x) dx (by Lemma 8.11) 0 Z n = 2xP (|Yn | > x) dx (since P (|Yn | > n) = 0) Z0 n = 2xP (|Xn | > x) dx (if |Yn | ≤ n then Yn = Xn ) Z0 n = 2xP (|X1 | > x) dx (identical distributions). 0

Next Z n ∞ ∞ X E(Yn2 ) X 1 = 2xP (|X1 | > x) dx n2 n2 0 n=1 n=1 Z ∞ ∞ X 1 = 2x1[0,n) (x)P (|X1 | > x) dx n2 0 n=1 Z ∞X ∞ 1 =2 x1[0,n) (x)P (|X1 | > x) dx. n2 0 n=1 P∞ We examine the function x 7→ n=1 n12 x1[0,n) (x). If 0 ≤ x ≤ 1 then none of the terms in the series is killed and ∞ ∞ X X 1 π2 1 x1 (x) ≤ = 1, then we only have the sum over n ≥ x. Let m = Int(x) (the integer part of x). We have Z ∞ ∞ ∞ X X 1 1 1 1 x1 (x) = x ≤ x dx = x ≤ 2. [0,n) 2 2 2 n n x m m n=1 n=m+1 In each case the function in question is bounded by 2 so Z ∞ ∞ X E(Yn2 ) ≤ 4 P (|X1 | > x) dx = 4E(|X1 |) < ∞. n2 0 n=1

268


Consider Yn − EYn and apply the previous theorem to get n

1X (Yk − EYk ) → 0 almost surely. n k=1

We have

E(Yk ) = E(Xk 1{ω:|Xk |≤k} ) = E(X1 1{ω:|X1 (ω)|≤k} ) → m since X1 1{ω:|X1 (ω)|≤k} → X1 , the sequence being dominated by |X1 | which is integrable. So we have by the triangle inequality |

n

n

n

k=1

k=1

k=1

1X 1X 1X Yk − m| ≤ | (Yk − EYk )| + |EYk − m| → 0. n n n

As observed earlier, this implies almost sure convergence of

Sn n .

8.2.5 Weak convergence In order to derive central limit theorems we first need to investigate the convergence of the distributions of the sequence (Xn ) of random variables.

Definition 8.4 A sequence Pn of Borel probability measures on Rn converges weakly to P if and only their cumulative distribution functions Fn converge to the distribution function F of P at all points where F is continuous. If Pn = PXn , P = PX are the distributions of some random variables, then we say that the sequence (Xn ) converges weakly to X. The name ‘weak’ is justified since this convergence is implied by the weakest we have come across so far, i.e. convergence in probability.

Theorem 8.22 If Xn converge in probability to X, then the distributions of Xn converge weakly.

Proof Let F (y) = P (X ≤ y). Fix y, a continuity point of F , and ε > 0. The goal is to obtain F (y) − ε < Fn (y) < F (y) + ε

8. Limit theorems

269

for sufficiently large n. By continuity of F we can find δ > 0 such that P (X ≤ y) −

ε < P (X ≤ y − δ), 2

ε P (X ≤ y + δ) < F (y) + . 2

By convergence in probability, P (|Xn − X| > δ)
X + (ω). Then F (y) > ω and, by the weak convergence, for sufficiently large n we have Fn (y) > ω. Then, by the construction, Xn+ (ω) ≤ y. This inequality holds for all except finitely many n so it is preserved if we take the upper limit on the left: lim sup Xn+ (ω) ≤ y.

270


Take a sequence yk of continuity points of F converging to X + (ω) from above (the set of discontinuity points of a monotone function is at most countable). For y = yk consider the above inequality and pass to the limit with k to get lim sup Xn+ (ω) ≤ X + (ω). Similarly lim inf Xn− (ω) ≥ X − (ω) so X − (ω) ≤ lim inf Xn− (ω) ≤ lim sup Xn+ (ω) ≤ X + (ω). The extremes are equal a.s. so the convergence holds a.s. The Skorokhod theorem is an important tool in probability. We will only need it for the following result, which links convergence of distributions to that of the associated characteristic functions.

Theorem 8.24 If PXn converge weakly to PX then ϕXn → ϕX .

Proof Take the Skorokhod representation Yn , Y of the measures PXn , PX . Almost sure convergence of Yn to Y implies that E(eitYn ) → E(eitY ) by the dominated convergence theorem. But the distributions of Xn , X are the same as the distributions of Yn , Y , so the characteristic functions are the same.

Theorem 8.25 (Helly’s Theorem) Let Fn be a sequence of distribution functions of some probability measures. There exists F , the distribution function of a measure (not necessarily probability), and a sequence kn such that Fkn (x) → F (x) at the continuity points of F.

Proof Arrange the rational numbers in a sequence: Q = {q1 , q2, . . .}. The sequence Fn (q1 ) is bounded (the values of a distribution function lie in [0,1]), hence it has a convergent subsequence, Fkn1 (q1 ) → y1 .

8. Limit theorems

271

Next consider the sequence Fkn1 (q2 ), which is again bounded, so for a subsequence kn2 of kn1 we have convergence Fkn2 (q2 ) → y2 . Of course also Fkn2 (q1 ) → y1 .

Proceeding in this way we find kn3 , kn4 ,. . . , each term a subsequence of the previous one, with Fkn3 (qm ) → ym for m ≤ 3, Fkn4 (qm ) → ym for m ≤ 4,

and so on. The diagonal sequence Fkn = Fknn converges at all rational points. We define FQ on Q by FQ (q) = lim Fkn (q) and next we write F (x) = inf{FQ (q) : q ∈ Q, q > x}.

We show that F is non-decreasing. Since Fn are non-decreasing, the same is true for FQ (q1 < q2 implies Fkn (q1 ) ≤ Fkn (q2 ) which remains true in the limit). Now let x1 < x2 . F (x1 ) ≤ FQ (q) for all q > x1 hence in particular for all q > x2 , so F (x1 ) ≤ inf q>x2 FQ (q) = F (x2 ). We show that F is right-continuous. Let xn & x. By the monotonicity of F , F (x) ≤ F (xn ) hence F (x) ≤ lim F (xn ). Suppose that F (x) < lim F (xn ). By the definition of F there is q ∈ Q, x < q, such that FQ (q) < lim F (xn ). For some n0 , x ≤ xn0 < q hence F (xn0 ) ≤ FQ (q) again by the definition of F , thus F (xn0 ) < lim F (xn ) which is a contradiction. Finally, we show that if F is continuous at x, then Fkn (x) → F (x). Let ε > 0 be arbitrary and find rationals q1 < q2 < x < q3 such that F (x) − ε < F (q1 ) ≤ F (x) ≤ F (q3 ) < F (x) + ε. Since Fkn (q2 ) → FQ (q2 ) ≥ F (q1 ), for sufficiently large n F (x) − ε < Fkn (q2 ). But Fkn is non-decreasing, so Fkn (q2 ) ≤ Fkn (x) ≤ Fkn (q3 ). Finally, Fkn (q3 ) → FQ (q3 ) ≥ F (q3 ), so for sufficiently large n Fkn (q3 ) < F (x) + ε. Putting together the above three inequalities we get F (x) − ε < Fkn (x) < F (x) + ε which proves the convergence.

272


Remark 8.4 The limit distribution function need not correspond to a probability measure. Example: Fn = 1[n,∞) , Fn → 0 so F = 0. This is a distribution function (nondecreasing, right continuous) and the corresponding measure satisfies P (A) = 0 for all A. We then say informally that the mass escapes to infinity. The following concept is introduced to prevent this happening.

Definition 8.5 We say that a sequence of probabilities Pn on Rd is tight if for each ε > 0 there is M such that Pn (Rd \ [−M, M ]) < ε for all n. By an interval in Rn we understand the product of intervals: [−M, M ] = {x = (x1 , . . . , xn ) ∈ Rn : xi ∈ [−M, M ] all i}. It is important that the M chosen for ε is good for all n – the inequality is uniform. It is easy to find such an M = Mn for each Pn separately. This follows from the fact that Pn ([−M, M ]) → 1 as M → ∞.

Theorem 8.26 (Prokhorov’s Theorem) If a sequence Pn is tight, then it has a subsequence convergent weakly to some probability measure P .

Proof By Helly’s theorem a subsequence Fkn converges to some distribution function F . All we have to do is to show that F corresponds to some probability measure P , which means we have to show that F (∞) = 1 (i.e. limy→∞ F (y) = 1). Fix ε > 0 and find a continuity point such that Fn (y) = Pn ((−∞, y]) > 1 − ε for all n (find M from the definition of tightness and take a continuity point of F which is larger than M ). Hence limn→∞ Fkn (y) ≥ 1 − ε, but this limit is F (y). This proves that limy→∞ F (y) = 1. We need to extend the notion of the characteristic function.

Definition 8.6 We sayR that ϕ is the characteristic function of a Borel measure P on R if ϕ(t) = eitx dP (x).

8. Limit theorems

273

In the case where P = PX we obviously have ϕP = ϕX , so the two definitions are consistent.

Theorem 8.27 Suppose (Pn ) is tight and let P be the limit of a subsequence of (Pn ) as provided by Prokhorov’s theorem. If ϕn (u) → ϕ(u) where ϕn are the characteristic functions of Pn and ϕ is the characteristic function of P , then Pn → P weakly.

Proof Fix a continuity point of F . For every subsequence Fkn there is a subsequence (subsubsequence) lkn , ln for brevity, such that Fln converge to some function H (Helly’s theorem). Denote the corresponding measure by P 0 . Hence ϕln → ϕP 0 , but on the other hand, ϕln → ϕ. So ϕP 0 = ϕ and consequently P 0 = P (Corollary 6.18). The above is sufficient for the convergence of the sequence Fn (y).

8.2.6 Central Limit Theorem The following lemma will be useful in what follows. It shows how to estimate the ‘weight’ of the ‘tails’ of a probability measure, and will be a useful tool in proving tightness.

Lemma 8.28 If ϕ is the characteristic function of P , then P (R \ [−M, M ]) ≤ 7M

Z

1 M

0

[1 − 1 (log m → ∞ as m → ∞) – the sequence is not Cauchy.

9. Solutions to exercises

293

Rm 1 1 (c) kfn − fm k1 = n x12 dx = − x1 |m n = n − m (n < m as before), and 1 ε 1 for any ε > 0 take N such that N < 2 and for n, m ≥ N , clearly n1 − m m so the sequence is not Cauchy. m 1 (b) kfn − fm k22 = R n x12 dx = n1 − m → 0 – the sequence is Cauchy. m 1 1 1 2 (c) kfn − fm k2 = n x4 dx = ( 3n3 − 3m 3 ) – the sequence is Cauchy.

5.5 kf + gk2 = 4, kf − gk2 = 1, kf k2 = 1, kgk2 = 1, and the parallelogram law is violated. 5.6 kf + gk21 = 0, kf − gk21 = parallelogram law.

1 2,

kf k21 =

1 4,

kgk21 =

1 4,

which contradicts the

5.7 Since sin nx cos mx = 12 [sin(n + m)x + sin(n − m)x] and sin nx sin mx = 1 2 [cos(n − m)x + cos(n + m)x], it is easy to compute the indefinite integrals. They are periodic functions so the integrals over [−π, π] are zero (for the latter we need n 6= m).

5.8 No: take any n, m (suppose n < m) and compute 1/n 4 Z n1 1 √ kgn − gm k44 = dx = −x−1 = (m − n) ≥ 1 1 x m 1/m

so the sequence is not Cauchy.

R1 5.9 Let Ω = [0, 1] with Lebesgue measure, X(ω) = √1ω , E(X) = 0 X dm = R1 1 R1 1 √ dx = 2, E(X 2 ) = √1 − 2 then 0 0 x dx = ∞. If we take X(ω) = x ω 2 E(X) = 0 and E(X ) = ∞.

5.10 Var(aX) = E((aX)2 ) − (E(aX))2 = a2 (E(X 2 ) − (E(X))2 ) = a2 Var(X). 2

(a+b) 1 2 1[a,b] , E(X) = a+b , E(X 2 ) = 5.11 Let fX (x) = b−a 2 , VarX = E(X ) − 4 R b 1 1 1 3 2 3 b−a a x dx = b−a 3 (b − a ) and simple algebra gives the result.

5.12 (a) E(X) = a, E((X − a)2 ) = 0 since X = a a.s. (b) By Exercise 4.15 fX (x) = 21[0, 12 ] (x) and by Exercise 4.17, E(X) = 14 ; R1 R1 1 so Var(X) = 02 2(x − 14 )2 dx = 2 −4 1 x2 dx = 48 . 4

(c) By Exercise 4.15, fX = 4x1[0, 12 ] (x), and by Exercise 4.17, E(X) = 16 , R1 1 hence Var(X) = 4 02 x(x − 61 )2 dx = 48 .

5.13 Cov(Y, 2Y +1) = E((Y )(2Y +1))−E(Y )E(2Y +1) = 2E(Y 2 )−2(E(Y ))2 = 2Var(Y ), Var(2Y + 1) = Var(2Y ) = 4Var(Y ), hence ρ = 1. 5.14 X, Y are uncorrelated by Exercise 5.7. Take a > 0 so small that the sets A = {ω : sin 2πω > 1 − a}, B = {ω : cos 2πω > 1 − a} are disjoint. Then P (A ∩ B) = 0 but P (A)P (B) 6= 0.

294


Chapter 6 6.1 The function g(x, y) =

 

1 x2

− 12  y 0

if 0 < y < x < 1 if 0 < x < y < 1 otherwise

is not integrable since the integral of g + is infinite (the same is true for the integral of g − ). However, Z 1Z 1 Z 1Z 1 g(x, y) dx dy = −1, g(x, y) dy dx = 1 0

0

0

0

which shows that the iterated integrals may not be equal if Fubini’s theorem condition is violated. RR R2 R3 2 R2 27 2 6.2 [0,3]×[−1,2] x y dm2 = −1 0 x y dx dy = −1 9y dy = 2 .

6.3 By Rsymmetry it is sufficient to consider x ≥ 0, y ≥ 0, and hence the area is a√ 4 ab 0 a2 − x2 dx = abπ. R2 6.4 Fix x ∈ [0, 2], 0 1A (x, y) dy = m(Ax ), hence fX (x) = x for x ∈ [0, 1], fX (x) = 2 − x for x ∈ (1, 2] and zero otherwise (triangle distribution). By symmetry, the same holds for fY . RR 6.5 P (X + Y > 4) = P (Y > −X + 4) = A fX,Y (x, y) dxRdyR where A = 2 4 1 {(x, y) : y > 4 − x} ∩ [0, 2] × [1, 4], so P (X + Y > 4) = 0 4−x 50 (x2 + R 2 8 1 (−4x2 + 43 x3 + 16x)) dx = 15 . y 2 ) dy dx = 50 0 P (Y > X) =

Z 4Z 2 1 2 1 2 (x + y 2 ) dx dy + (x + y 2 ) dx dy 50 50 1 0 2 0 Z 4 4 3 1 8 2 y dy + ( + 2y ) dy 3 50 2 3

Z

Z 2 1 = 50 1 143 = 150

2

Z

y

RR Similarly P (Y > X) = fX,Y (x, y) dx dy where A = {(x, y) : y > x} ∩ R 2 RAy 1 2 R4R2 1 2 [0, 2] × [1, 4], so we get 1 0 50 (x + y 2 ) dx dy + 2 0 50 (x + y 2 ) dy dx = R R 2 4 3 4 8 1 1 143 2 50 1 3 y dy + 50 2 ( 3 + 2y ) dy = 150 .

6.6

 0    z fX+Y (z) = fX,Y (x, z − x) dx =  2−z R   0 Z

z 0) = 0 fY −X (z) dz = 21 + 1 12 (2 − z) dz = 1 1 3 2 + 4 = 4; R∞ R3 P (X + Y > 1) = 1 fX+Y (z) dz = 12 + 2 12 (3 − z) dz = 12 + 13 = 43 . R − 1 x+1 A (x,y) 6.9 fX (x) = 0 2 1A dy = 1 − 21 x, h(y, x) = 11− and E(Y |X = 1) = 1 2x R 12 2 0 x dx = 41 . R1 6.10 fY (y) = 0 (x + y) dx = 12 + y, h(x|y) = x+y 1A (x, y), E(X|Y = y) = 1 2 +y 1 1 R 1 x+y + y 2 . x 1 +y dx = 31 +y 0 2

2

Chapter 7 7.1 If µ(A) = 0 then λ1 (A) = 0 and λ2 (A) = 0, hence (λ1 + λ2 )(A) = 0. 7.2 Let Q be a finite partition of Ω which refines both P1 and P2 . Thus each Sn Sm set A ∈ Q can be written as a disjoint union A = i=1 Ei = j=1 Fj where Ei ∈ P1 , Fj ∈ P2 . Each element of A belongs to exactly one Ei and exactly S one Fj so A = i,j (Ei ∩ Fj ) is a disjoint union as well. Hence Q refines the partition R = {E ∩ F : E ∈ P1 , F ∈ P2 } (which is a partition as the above argument applied to A = Ω shows). It is sufficient to see that R refines Pi , S S i = 1, 2. But E ∈ P1 can be written as E = E ∩ F ∈P2 F = F ∈P2 (E ∩ F ) so E is a disjoint union of elements of R. Similar argument shows that R refines P2 .

7.3 We have to assume first that m(A) 6= 0. Then B ⊂ A clearly implies that µ dominates ν. (In fact m(B \ A) = 0 is slightly more general.) Then consider the partitionR{B, A\B, Ω \A} R to see that h = 1B . To double check, ν(F ) = m(F ∩ B) = F ∩B 1B dm = F ∩B 1B dµ.

296


7.4 Clearly µ({ω}) ≥ ν({ω}) is equivalent to µ dominating ν. For each ω we ν({ω}) dν (ω) = µ({ω}) . have dµ R R dν R dν 7.5 Since ν(E) = E g dm and we wish to have ν(E) = E dµ dµ = E dµ f dm

dν it is natural to aim at taking dµ (x) = fg(x) (x) . Then a sufficient condition for this to work is that if A = {x : f (x) = 0} then ν(A) = 0, i.e. g(x) = 0 dν c a.e. on A. Then we put dµ (x) = fg(x) (x)R on A and 0 on A and we have R R dν dν ν(E) = E∩Ac g dm = E∩Ac dµ f dm = E dµ dµ, as required.

7.6 Clearly ν µ is equivalent to A = {ω : µ({ω}) = 0} ⊂ {ω : ν({ω}) = 0} ν({ω}) dν (ω) = µ({ω}) on A and zero outside. and then dµ R dν 7.7 Since ν µ we may write h = dµ , so that ν(F ) = F h dµ. As µ(F ) = 0 if and only if ν(F ) = 0, the set {h = 0} is both µ-null and ν-null. Thus dν −1 ) is well-defined a.s., and we can use (ii) in Proposition 7.7 h−1 = ( dµ with λ = µ to conclude that 1 = h−1 h implies

dµ dν

= h−1 , as required.

1 1 7.8 δ0 ((0, 25]) = 0, but 25 m|[0,25] ((0, 25]) = 1; 25 m|[0,25] ({0}) = 0 but δ0 ({0}) = dP1 (x) = 1 so neither P1 P2 nor P2 P1 . Clearly P1 P3 with dP 3 dP2 2 × 1{0} (x) and P2 P3 with dP3 (x) = 2 × 1(0,25] (x).

7.9 λa = m|[2,3] , λs = δ0 + m|(1,2) , and h = 1[2,3] .

7.10 Suppose F is non-constant at ai with positive jumps ci , i = 1, 2, . . . Take M 6= ai , with −M 6= ai and let I = {i : ai ∈ [−M, M ]}. Then X X ci = mF ({ai }), mF ([−M, M ]) = F (M ) − F (−M ) = i∈I

i∈I

which is finite since F is bounded on a bounded interval. So any A ⊂ S [−M, M ] \ i∈I {ai } is mF -null hence measurable. But {ai } are mF measurable hence each subset of [−M, M ] is mF -measurable. Finally, any subset E of R is a union of the sets of the form E ∩ [−M, M ], so E is mF -measurable as well. 7.11 mF has density f (x) = 2 for x ∈ [0, 1] and zero otherwise. Rx 7.12 (a) |x| = 1 + −1 f (y) dy, where f (y) = −1 for y ∈ [−1, 0], and f (y) = 1 for y ∈ (0, 1]. Pn (b) Let 1 > ε > 0, take δ = ε2 , k=1 (yk − xk ) < δ, with yk ≤ xk+1 ; then n X √ √ √ √ √ ( | xk − yk |)2 ≤ ( yn − x1 )2 = yn − 2 yn x1 + x1 < yn − x1 k=1

≤

n X

k=1

(yk − xk ) < ε2 .


297

(c) Lebesgue function f is a.e. differentiable with f 0 = R x0 0a.e. If it were absolutely continuous, it could be written as f (x) = 0 f (y) dy = 0, a contradiction. Pk 7.13 (a) If F is monotone increasing on [a, b], i=1 |F (xi ) − F (xi−1 )| = F (b) − F (a) for any partition a = x0 < x1 < · · · < xk = b. Hence TF [a, b] = F (b) − F (a). (b) If F ∈ BV [a, b] we can write F = F1 − F2 where both F1 , F2 are monotone increasing, hence have only countably many points of discontinuity. So F is continuous a.e. and thus Lebesgue-measurable. (c) f (x) = x2 cos xπ2 for x 6= 0 and f (0) = 0 is differentiable but does not belong to BV [0, 1]. Pk Pk (d) If F is Lipschitz, i=1 |F (xi ) − F (xi−1 )| ≤ M i=1 |xi − xi−1 | = M (b − a) for any partition so TF [a, b] ≤ M (b − a) is finite.

7.14 Recall that ν + (E) = ν(E ∩ B), where B is the positive set in the Hahn decomposition. As in the hint, if G ⊆ F, ν(G) ≤ ν + (G ∩ B) ≤ ν(G ∩ B) + ν((F ∩ B) \ (G ∩ B)) = ν(F ∩ B). Since the set (F ∩ B) \ (G ∩ B) ⊆ B, its ν-measure is non-negative. But F ∩ B ⊆ F so sup{ν(G) : G ⊆ F } is attained and equals ν(F ∩ B) = ν + (F ). A similar argument shows ν − (F ) = sup{−ν(G)} = − inf G⊂F {ν(G)}. R R 7.15 For all F ∈ F,ν + (F ) = B∩F fR dµ = supG⊂F G f Rdµ. If f >R 0 on a set C ⊂ A ∩RF with µ(C) > 0, then C f dµ > 0, so that C f dµ + B∩F f dµ > supG⊂F G f dµ. This is a contradiction since C ∪ (B ∩ F ) ⊂ F. So f ≤ 0 a.s. (µ) on A∩F . We can take the set {f = 0} into B, since it does not affect the integrals. Hence {f < 0} ⊂ A and {f ≥ 0} ⊂ B. But the two smaller sets partition Ω, so we have equality in both cases. Hence f + = f 1B and f − = −f 1A , therefore for all F ∈ F Z Z ν + (F ) = ν(B ∩ F ) = f dµ = f + dµ, B∩F F Z Z − ν (F ) = −ν(A ∩ F ) = − f dµ = f − dµ. A∩F

1

+

F

R 7.16 fR ∈ L (ν) iff both f dν and f dν are finite. Then ERf + g dµ and − well-defined and finite and their difference is E f g dµ. So E f g dµ are R fR g ∈ L1 (µ), as ER(f + −f − )|g| dµ < ∞. Conversely, if f g ∈ L1 (µ) then R both + − f |g| dµ and f |g| dµ are finite hence so is their difference E f dν. E E R

R

−

7.17 (a) E(X|G)(ω) = 14 if ω ∈ [0, 12 ], E(X|G)(ω) = 43 otherwise. (b) E(X|G)(ω) = ω if ω ∈ [0, 12 ], E(X|G)(ω) = 43 otherwise.

7.18 E(Xn |Fn−1 ) = E(Z1 Z2 . . . Zn |Fn−1 ) = Z1 Z2 . . . Zn−1 E(Zn |Fn−1 ), and since Zn is independent of Fn−1 , E(Zn |Fn−1 ) = E(Zn ) = 1, hence the result.

298


7.19 E(Xn ) = nµ 6= µ = E(X1 ) so Xn is not a martingale. Clearly Yn = Xn − nµ is a martingale. 7.20 E((Z1 + Z2 + · · · + Zn )2 |Fn−1 ) ≤ E(Z1 + Z2 + · · · + Zn |Fn−1 )2 = (Z1 + Z2 +· · ·+Zn−1 )2 using Jensen inequality. The compensator is deterministic: An = n. 7.21 For s < t, since the increments are independent and w(t) − w(s) has the same distribution as w(t − s), 1 1 2 E(exp(−σw(t) − α2 t)) = e−σw(s)− 2 σ t E(exp(−[σ(w(t) − w(s)])|Fs ) 2 1 2 = e−σw(s)− 2 σ t E(exp(−[σ(w(t) − w(s)]) 1

2

= e−σw(s)− 2 σ t E(exp(−σw(t − s))).

Now σw(t − s) v N (0, σ 2 (t − s)) so the expectation equals E(e−σ 1 2 e− 2 σ (t−s) (where Z v N (0, 1)) and so the result follows.

√ t−sZ

)=

Chapter 8 8.1 (a) fn = 1[n,n+ n1 ] converges to 0 in Lp , pointwise, a.e. but not uniformly. (b) fn = n1[0, n1 ] −n1[− n1 ,0] converges to 0 pointwise and a.e. It converges neither in Lp nor uniformly. 8.2 We have Ω = [0, 1] with Lebesgue measure. The sequences Xn = 1(0, n1 ) , Xn = n1(0, n1 ] converge to 0 in probability since P (|Xn | > ε) ≤ n1 and the same holds for the sequence gn . 8.3 There are endless possibilities, the simplest being Xn (ω) ≡ 1 (but this sequence converges to 1) or, to make sure that it does not converge to anything, Xn (ω) ≡ n.

100 8.4 Let Xn = 1 indicate the heads and Xn = 0 the tail, then S100 is the S100 1 average number of heads in 100 tosses. Clearly E(Xn ) = 2 , E( 100 ) = 21 , 1 1 1 100 Var(Xn ) = 41 , Var( S100 ) = 100 2 100 · 4 = 400 so

P (|

S100 1 1 − | ≥ 0.1) ≤ 100 2 0.12 400

and

S100 1 1 3 − | < 0.1) ≥ 1 − = . 100 2 0.12 400 4 8.5 Let Xn be the number shown on the die, E(Xn ) = 3.5, Var(Xn ) ≈ 2.9. P (|

P (|

S1000 − 3.5| < 0.01) ≥ 0.29. 1000


8.6 The union

299

S∞

Am is equal to [0, 1] for all m and so is lim supn→∞ An . 1 8.7 Let d = 1. There are 2n paths that return to 0, so P (S2n = 0) = 2n n n 22n . Now √ √ 2n ( 2n 2π2n (2n)! 2n 2 e ) ∼ = √ n 2n 2 (n!) ( e ) 2πn nπ q P∞ so P (S2n = 0) ∼ √cn with c = π2 . Hence n=1 P (An ) diverges and BorelCantelli applies (as (An ) are independent) so that P (S2n = 0 i.o.) = 1. 1 Same for d = 2 since P (An ) ∼ n1 . But for d > 2, P (An ) ∼ nd/2 , the series converges and by the first Borel-Cantelli lemma P (S2n = 0 i.o.) = 0. m=n

√ 8.8 Write S = S1000 ; P (|S − 500| < 10) = P ( |S−500| < 0.63) ≈ 0.47. 250

√ 8.9 The condition on n is P (| Snn − 0.5| < 0.005) = P ( |Sn√−0.5n| < 0.01 n) ≥ n/4

0.99, hence n ≥ 66 615. √ 8.10 Write xn = eσ T /n . Then

T 1 1 (2Rn xn )2 1 + x2n 1 (ln Un + ln Dn ) = ln(Un Dn ) = ln = ln er n n − ln( ). 2 2 2 2 2 (1 + xn ) 2xn 2

So it suffices to show that the last term on the right is σ2nT + o( n1 ). But √ √ p x−1 eσ T /n + eσ T /n 1 + x2n n + xn = = = cosh(σ T /n) 2xn 2 2 σ2 T 1 =1+ + o( ) 2n n so that ln(

1 + x2n σ2 T 1 σ2 T 1 ) = ln(1 + + o( )) = + o( ). 2xn 2n n 2n n

10 Appendix

Existence of non-measurable and non-Borel sets In Chapter 2 we defined the σ-field B of Borel sets and the larger σ-field M of Lebesgue-measurable sets, and all our subsequent analysis of the Lebesgue integral and its properties involved these two families of subsets of R. The set inclusions B ⊂ M ⊂ P(R) are trivial; however, it is not at all obvious at first sight that they are strict, i.e. that there are sets in R which are not Lebesgue-measurable, as well as that there are Lebesgue-measurable sets which are not Borel sets. In this appendix we construct examples of such sets. Using the fact that A ⊂ R is measurable (resp. Borel-measurable) iff its indicator function 1A ∈ M (resp. B) it follows that we will automatically have examples of non-measurable (resp. measurable but not Borel) functions. The construction of a non-measurable set requires some set-theoretic preparation. This takes the form of an axiom which, while not needed for the consistent development of set theory, nevertheless enriches that theory considerably. Its truth or falsehood cannot be proved from the standard axioms on which modern set theory is based, but we shall accept its validity as an axiom, without delving further into foundational matters.

301

302


The Axiom of Choice Suppose that A = {Aα : α ∈ Λ} is a non-empty collection, indexed by some set Λ, of non-empty disjoint subsets of a fixed set Ω, Then there exists a set E ⊂ Ω which contains precisely one element from each of the sets Aα , i.e. there is a choice function f : Λ → A.

Remark The Axiom may seem innocuous enough, yet it can be shown to be independent of the (Zermelo–Fraenkel) axioms of sets theory. If the collection A has only finitely many members there no problem in finding a choice function, of course. To see that the existence of such a function is problematic for infinite sets, consider the following illustration given by Bertrand Russell: imagine being faced with an infinite collection of pairs of shoes and another of pairs of socks. Constructing the set consisting of all left shoes is simple; that of defining the set of all left socks is not! To construct our example of a non-measurable set, first define the following equivalence relation on [0, 1]: x ∼ y if y − x is a rational number (which will be in [−1, 1]). This relation is easily seen to be reflexive, symmetric and transitive. Hence it partitions [0, 1] into disjoint equivalence classes (Aα ), where for each α, any two elements x, y of Aα differ by a rational, while elements of different classes will always differ by an irrational. Thus each Aα is countable, since Q is, but there are uncountably many different classes, as [0, 1] is uncountable. Now use the Axiom of Choice to construct a new set E ⊂ [0, 1] which contains exactly one member aα from each of the Aα . Now enumerate the rationals in [−1, 1]: there are only countably many, so we can order them as a sequence (qn ). Define a sequence of translates of E by En = E + qn . If E is Lebesgue-measurable, then so is each En and their measures are the same, by Proposition 2.10. But the (En ) are disjoint: to see this, suppose that z ∈ Em ∩ En for some m 6= m. Then we can write aα + qm = z = aβ + qn for some aα , aβ ∈ E, and their difference aα −aβ = qn −qm is rational. Since E contains only one element S∞ from each class, α = β and therefore m = n. Thus n=1 En is a disjoint union containing [0, 1]. S∞ Thus we have [0, 1] ⊂ n=1 En ⊂ [−1, 2] and m(En ) = m(E) for all n. By countable additivity and monotonicity of m this implies: 1 = m([0, 1]) ≤

∞ X

n=1

m(En ) = m(E) + m(E) + · · · ≤ 3.

This is clearly impossible, since the sum must be either 0 or ∞. Hence we must conclude that E is not measurable.

10. Appendix

303

For an example of a measurable set that is not Borel, let C denote the Cantor set, and define the Cantor function f : [0, 1] → C as follows: for x ∈ [0, 1] write P an x = 0.a1 a2 . . . in binary form, i.e. x = ∞ n=1 2n , where each an = 0 or 1 (taking non-terminating expansions where the choice exists). The function x 7→ an is determined by a system of finitely many binary intervals (i.e. the value of an is fixed by x satisfying finitely many linear inequalities) and so is measurable P∞ n – hence so is the function f given by f (x) = n=1 2a 3n . Since all the terms of P∞ 2an y = n=1 3n have numerators 0 or 2, it follows that the range Rf of f is a subset of C. Moreover, the value of y determines the sequence (an ) and hence x, uniquely, so that f is invertible. Now consider the image in C of the non-measurable set E constructed above, i.e. let B = f (E). Then B is a subset of the null set C, hence by the completeness of m it is also measurable and null. On the other hand, E = f −1 (B) is non-measurable. We show that this situation is incompatible with B being a Borel set. Given a set B ∈ B and a measurable function g, then g −1 (B) must be measurable. For, by definition of measurable functions, g −1 (I) is measurable for every interval I, and we have g −1 (

∞ [

i=1

Ai ) =

∞ [

g −1 (Ai ),

g −1 (Ac ) = (g −1 (A))c

i=1

quite generally for any sets and functions. Hence the collection of sets whose inverse images under the measurable function g are again measurable forms a σ-field containing the intervals, hence also contains all Borel sets. But we have found a measurable function f and a Lebesgue-measurable set B for which f −1 (B) = E is not measurable. Therefore the measurable set B cannot be a Borel set, i.e. the inclusion B ⊂ M is strict.

Bibliography

[1] T.M. Apostol, Mathematical Analysis, Addison–Wesley, Reading, 1974. [2] P.Billingsley, Probability and Measure, John Wiley and Sons, New York 1995 [3] Z.Brzezniak, T.Zastawniak, Basic Stochastic Processes, Springer–Verlag, London 1999 [4] M.Capinski, T.Zastawniak, Mathematics for Finance, An Introduction to Financial Engineering, Springer–Verlag, London 2003 [5] R.J.Elliott, P.E.Kopp, Mathematics of Financial markets, Springer–Verlag, New York 1999 [6] G.R. Grimmett and D.R. Stirzaker, Probability and Random Processes, Clarendon Press, Oxford, 1982. [7] J.Hull, Options, Futures, and Other Derivatives, Prentice Hall, Upper Saddle River NJ 2000 [8] P.E. Kopp, Analysis, Modular Mathematics, Edward Arnold, London, 1996. [9] J. Pitman, Probability, Springer–Verlag, New York, 1995. [10] W. Rudin, Real and Complex Analysis, McGraw–Hill, New York, 1966. [11] G. Smith, Introductory Mathematics: Algebra and Analysis, Springer– Verlag, SUMS, 1998. [12] D. Williams, Probability with Martingales, Cambridge University Press, Cambridge, 1991.

305

Index

a.e., 55 – convergence, 242 a.s., 56 absolutely continuous, 189 – function, 109, 204 – measure, 107 adapted, 222 additive – countably, 27 additivity – countable, 29 – finite, 39 – of measure, 35 almost everywhere, 55 almost surely, 56 American option, 72 angle, 138 Banach space, 136 Beppo-Levi – theorem, 95 Bernstein – polynomials, 250 Bernstein-Weierstrass – theorem, 250 binomial – tree, 50 Black-Scholes – formula, 118 – model, 118 Borel – function, 57 – measure, regular, 44 – set, 40 Borel-Cantelli lemma

– first, 257 – second, 258 bounded variation, 206 Brownian motion, 233 BV[a,b], 206 call option, 71 – down-and-out, 72 call-put parity, 117 Cantor – function, 303 – set, 19 Cauchy – density, 108 – sequence, 11, 128 central limit theorem, 276, 280 central moment, 146 centred – random variable, 151 characteristic function, 116, 272 Chebyshev’s inequality, 247 complete – measure space, 43 – space, 128 completion, 43 concentrated, 197 conditional – expectation, 153, 178, 179, 218 – probability, 47 contingent claim, 71, 72 continuity of measure, 39 continuous – absolutely, 204 convergence – almost everywhere, 242 307

308

– in Lp , 242 – in p-th mean, 144 – in probability, 245 – pointwise, 242 – uniform, 11, 241 – weak, 268 correlation, 138, 151 countable – additivity, 27, 29 covariance, 151 cover, 20 de Moivre–Laplace theorem, 280 de Morgan’s laws, 3 density, 107 – Cauchy, 108 – Gaussian, 107, 174 – joint, 173 – normal, 107, 174 – triangle, 107 derivative – Radon-Nikodym, 194 derivative security, 72 – European, 71 Dirac measure, 68 direct sum, 139 distance, 126 distribution – function, 109, 110, 199 – gamma, 109 – geometric, 69 – marginal, 174 – Poisson, 69 – triangle, 107 – uniform, 107 dominated convergence theorem, 92 dominating measure, 190 Doob decomposition, 226 essential – infimum, 66 – supremum, 66 essentially bounded, 141 event, 47 eventually, 256 exotic option, 72 expectation – conditional, 153, 178, 179, 218 – of random variable, 114 Fatou’s lemma, 82 filtration, 51, 222 – natural, 222 first hitting time, 230


formula – inversion, 180 Fourier series, 140 Fubini’s theorem, 171 function – Borel, 57 – Cantor, 303 – characteristic, 116, 272 – Dirichlet, 99 – essentially bounded, 141 – integrable, 86 – Lebesgue, 20 – Lebesgue measurable, 57 – simple, 76 – step, 102 fundamental theorem of calculus, 9, 97, 214 futures, 71 gamma distribution, 109 Gaussian density, 107, 174 geometric distribution, 69 H¨ older inequality, 142 Hahn-Jordan decomposition, 211, 216 Helly’s theorem, 270 Hilbert space, 136, 138 i.o., 255 identically distributed, 244 independent – events, 48, 49 – random variables, 70, 244 – σ-fields, 49 – σ-fields, 48 indicator function, 4, 59 inequality – Chebyshev, 247 – H¨ older, 142 – Jensen, 220 – Kolmogorov, 262 – Minkowski, 143 – Schwarz, 132, 143 – triangle, 126 infimum, 6 infinitely often, 255 inner product, 135, 136 – space, 136 integrable function, 86 integral – improper Riemann, 99 – Lebesgue, 77, 87 – of a simple function, 76 – Riemann, 7

Index

invariance – translation, 35 inversion formula, 180 Ito isometry, 229 Jensen inequality, 220 joint density, 173 Kolmogorov inequality, 262 L2 (E), 131 Lp (E), 140 L∞ (E), 141 law of large numbers – strong, 260, 266 – weak, 249 Lp convergence, 242 Lebesgue – decomposition, 197 – function, 20 – integral, 76, 87 – measurable set, 27 – measure, 35 Lebesgue-Stieltjes – measurable, 202 – measure, 199 lemma – Borel-Cantelli, 257 – Fatou, 82 – Riemann–Lebesgue, 104 Levy’s theorem, 274 liminf, 6 limsup, 6 Lindeberg–Feller theorem, 276 lower limit, 256 lower sum, 77 marginal distribution, 174 martingale, 223 – transform, 227 mean value theorem, 81 measurable – function, 57 – Lebesgue-Stieltjes, 202 – set, 27 – space, 189 measure, 29 – absolutely continuous, 107 – Dirac, 68 – F -outer, 200 – Lebesgue, 35 – Lebesgue-Stieltjes, 199 – outer, 20, 45 – probability, 46

309

– product, 164 – regular, 44 – σ-finite, 162 – signed, 209, 210 – space, 29 measures – mutually singular, 197 metric, 126 Minkowski inequality, 143 model – binomial, 50 – Black-Scholes, 118 – CRR, 233 moment, 146 monotone class, 165 – theorem, 165 monotone convergence theorem, 84 monotonicity – of integral, 81 – of measure, 21, 35 Monte-Carlo method, 251 mutually singular measures, 197 negative part, 63 negative variation, 207 norm, 126 normal density, 107, 174 null set, 16 option – American, 72 – European, 71 – exotic, 72 – lookback, 72 orthogonal, 137–139 orthonormal – basis, 140 – set, 139 outer measure, 20, 45 parallelogram law, 136 partition, 190 path, 50 pointwise convergence, 242 Poisson distribution, 69 polarization identity, 136 portfolio, 183 positive part, 63 positive variation, 207 power set, 2 predictable, 225 probability, 46 – conditional, 47 – distribution, 68

310

– measure, 46 – space, 46 probability space – filtered, 222 process – stochastic, 222 – stopped, 230 product – measure, 164 – σ-field, 160 Prokhorov’s theorem, 272 put option, 72 Radon-Nikodym – derivative, 194 – theorem, 190, 195 random time, 229 random variable, 66 – centred, 151 rectangle, 3 refinement, 7, 190 replication, 232 return, 183 Riemann – integral, 7 – – improper, 99 Riemann’s criterion, 8 Riemann–Lebesgue lemma, 104 Schwarz inequality, 132, 143 section, 162, 170 sequence – Cauchy, 11, 128 – tight, 272 set – Borel, 40 – Cantor, 19 – Lebesgue measurable, 27 – null, 16 σ-field, 29 – generated, 40 – product, 160 σ-field – generated – – by random variable, 67 σ-finite measure, 162 signed measure, 209, 210 simple function, 76 Skorokhod representation theorem, 110, 269 space – Banach, 136 – complete, 128 – Hilbert, 136, 138


– inner product, 136 – L2 (E), 131 – Lp (E), 140 – measurable, 189 – measure, 29 – probability, 46 standard normal distribution, 114 step function, 102 stochastic – integral, discrete, 227 – process, discrete, 222 stopped process, 230 stopping time, 229 strong law of large numbers, 266 subadditivity, 24 submartingale, 223 summable, 217 supermartingale, 224 supremum, 6 symmetric difference, 35 theorem – Beppo–Levi, 95 – Bernstein-Weierstrass Approximation, 250 – central limit, 276, 280 – de Moivre–Laplace, 280 – dominated convergence, 92 – Fubini, 171 – Fundamental of Calculus, 214 – fundamental of calculus, 9, 97 – Helly, 270 – intermediate value, 6 – Levy, 274 – Lindeberg–Feller, 276 – mean value, 81 – Miller-Modigliani, 117 – monotone class, 165 – monotone convergence, 84 – Prokhorov, 272 – Radon-Nikodym, 190, 195 – Skorokhod representation, 110, 269 tight sequence, 272 total variation, 207 translation invariance – of measure, 35 – of outer measure, 26 triangle inequality, 126 triangular array, 282 uncorrelated random variables, 151 uniform convergence, 11, 241 uniform distribution, 107 upper limit, 255

Index

upper sum, 77 variance, 147 variation – bounded, 206 – function, 207 – negative, 212

311

– positive, 212 – total, 207, 211 weak – convergence, 268 – law of large numbers, 249 Wiener process, 233

Measure, Integral and Probability

Measure, Integral and Probability

Measure, integral and probability