A
A15 phases — Anoxic zones
A15 phases A series of intermetallic compounds that have a particular crystal structure a...
345 downloads
2686 Views
18MB Size
Report
This content was uploaded by our users and we assume good faith they have the permission to share this book. If you own the copyright to this book and it is wrongfully on our website, we offer a simple DMCA procedure to remove your content from our site. Start by pressing the button below!
Report copyright / DMCA form
A
A15 phases — Anoxic zones
A15 phases A series of intermetallic compounds that have a particular crystal structure and the chemical formula A3B, where A represents a transition element and B can be either a transition element or a nontransition element. Many A15 compounds exhibit the phenomenon of superconductivity at relatively high temperatures in the neighborhood of 20 K (−424◦F) and in high magnetic fields on the order of several tens of teslas (several hundred kilogauss). High-temperature–high-field superconductivity has a number of important technological applications and is a challenging fundamental research area in condensed-matter physics. Crystal structure. The A15 compounds crystallize in a structure in which the unit cell, the repeating unit of the crystal structure, has the overall shape of a cube. The B atoms are located at the corners and in the center of the cube, while the A atoms are arranged in pairs on the cube faces (see illus.). A special characteristic of the A15 crystal structure is that the A atoms form mutually orthogonal linear chains that run throughout the crystal lattice, as shown in the illustration. The intrachain distance between A atoms is the shortest distance between atoms in the A15 crystal structure, and is about 22% less than the smallest interchain distance between A atoms. The extraordinary superconducting properties of the A15 compounds are believed to be primarily associated with these linear chains of transitionelement A atoms. At low temperature, some of the A15 phases undergo a diffusionless transformation from a cubic to a tetragonal structure, in which the length of one of the sides of the unit cell changes with respect to the length of the other two. This structural transformation was first observed in V3Si at approximately 21 K (−422◦F), and has also been found in Nb3Sn at 43 K (−382◦F), V3Ga at above 50 K (−370◦F), and Nb3Al at 80 K (−316◦F). See CRYSTAL STRUCTURE.
Superconductivity. Superconductivity is a phenomenon in which a metal, when cooled below its superconducting transition temperature Tc, loses all resistance to the flow of an electric current. The A15 compound Nb3Ge has a Tc of 23.2 K (−417.9◦F), the highest value of the A15 compounds, and, prior to 1986, the highest value ever observed for any known material. Other A15 compounds with high Tc include Nb3Al (18.8 K or −425.8◦F), Nb3Ga (20.3 K or −423.1◦F), Nb3Si (approximately 19 K or −426◦F), Nb3Sn (18.0 K or −427.3◦F), V3Ga (15.9 K or −431◦F), and V3Si (17.1 K or −428.9◦F). In 1986, evidence of superconductivity with a Tc of approximately 30 K (−405◦F) was reported for a material containing the elements lanthanum, barium, copper, and oxygen. Since then, several families of compounds containing copper, oxygen, and
A15 crystal structure of the A3B intermetallic compounds. The light spheres represent the A atoms; the dark spheres represent the B atoms. The linear chains of A atoms are emphasized.
2
Aardvark other elements have been discovered with yet higher superconducting transition temperatures, including a value of 135 K (−217◦F) for a compound made of mercury, thallium, barium, calcium, copper, and oxygen. The origin of high-temperature superconductivity in the copper oxide compounds has not yet been established and may be different from that in the conventional superconductors such as the A15 compounds. In a conventional superconductor, the interaction of the negatively charged electrons with the positively charged ions of the crystal lattice produces an attractive interaction between electrons, causing them to form electron pairs, called Cooper pairs, that are responsible for the superconductivity. The superconductivity of the A15 compounds can be destroyed with a sufficiently high magnetic field, called the upper critical magnetic field, Hc2(T), which depends on the temperature T and varies from material to material. As the temperature is lowered, Hc2(T) generally increases from zero at Tc to a maximum value, Hc2(0), as the temperature approaches absolute zero (0 K or −459.67◦F). The A15 compounds also have very high values of Hc2(0), the highest of which is 44 teslas (T) for a pseudobinary A15 compound with the composition Nb79(Al73Ge27)21. Other A15 compounds with high values of Hc2(0) include Nb3Al (32 T), Nb3Ge (39 T), Nb3Sn (23 T), V3Ga (21 T), and V3Si (25 T). For comparison, the highest value of Hc2(0) ever observed for a conventional superconducting material is about 60 T for the Chevrel-phase compound PbMo6S8, while the values of Hc2(0) for the highest-Tc copper oxide superconductors may reach several hundred teslas, values so large that they cannot readily be measured by using the magnetic fields and techniques that are presently available in the laboratory. See CHEVREL PHASES; SUPERCONDUCTIVITY. Technological applications. Because of their high values of Tc and Hc2, and critical current density Jc, the A15 compounds have a number of important technological applications. (The critical current density is the highest electric current per unit area that a material can conduct and still remain superconducting in a given magnetic field.) Processes have been developed for preparing multifilamentary superconducting wires that consist of numerous filaments of a superconducting A15 compound, such as Nb3Sn, embedded in a nonsuperconducting copper matrix. Superconducting wires can be used in electric power transmission lines and to wind electrically lossless coils (solenoids) for superconducting electrical machinery (motors and generators) and magnets. Superconducting magnets are employed to produce intense magnetic fields for laboratory research, confinement of high-temperature plasmas in nuclear fusion research, bending beams of charged particles in accelerators, levitation of high-speed trains, mineral separation, nuclear magnetic resonance imaging, and energy storage. If conductors of the hightemperature, high-field oxide superconductors can eventually be fabricated with high values of Jc, they may revolutionize the technological applications
of superconductivity. See SUPERCONDUCTING DEVICES. M. Brian Maple Bibliography. W. Buckle, Superconductivity: Fundamentals and Applications, 1991; P. Fulde (ed.), Superconductivity of Transition Metals, Springer Series in Solid State Sciences, vol. 27, 1982; T. H. Geballe and J. K. Hulm, Superconductivity: The state that came in from the cold, Science, 239:367– 374, 1988; Y. Kao (ed.), Superconductivity and Its Applications, 1992; V. Z. Kresin, H. Morawitz, and S. A. Wolf, Mechanics of Conventional and High Tc Superconductivity, 1993.
Aardvark A large piglike mammal in the order Tubulidentata. The order contains a single family, Orycteropodidae, and a single species, Orycteropus afer. Aardvarks occur throughout Africa south of the Sahara, wherever suitable habitat exists. In the past, tubulidentates were often considered closely related to ungulates. However, recent data from mitochondrial and nuclear genes support a relationship among aardvarks, elephant-shrews, paenungulates (hyraxes, sirenians, and proboscideans), and golden moles (Chrysochloridae). These ecologically divergent adaptive types probably originated in Africa; the molecular evidence implies that they may have arisen there from a common ancestor that existed in the Cretaceous Period, when Africa was isolated from other continents. This order shows the results of an extreme adaptation for burrowing and feeding on small food items (mainly termites). See MAMMALIA; TUBULIDENTATA. Characteristics. Aardvarks, also known as antbears, resemble a medium-sized to large pig (see illustration). The body is massive with an elongate head and snout. The round, blunt, piglike muzzle ends in circular nostrils that possess fleshy tentacles on the nasal septum and dense tracts of white hairs that can serve as a filter and seal the nostrils to prevent soil from entering the lungs. The tough, thick skin is sparsely covered with bristly hair ranging from dull brownish gray to dull yellowish
Aardvark, Orycteropus afer. (Photo by Lloyd Glenn Ingles; c 2001 California Academy of Sciences)
Abaca gray. Numerous vibrissae (hairs with specialized erectile tissues) are present on the face. The waxy, smooth, tubular donkeylike ears are large and can be folded back to exclude dirt when the animal is burrowing, and they can be moved independently of each other. The strong, muscular tail is kangaroolike, tapering to a point. Incisors and canines are lacking in the dentition of adults but are present in the deciduous dentition. Teeth on the sides of the jaw are simple, peglike, open-rooted (not anchored in the jaw), and grow continuously during the animal’s lifetime. The teeth are hollow; instead of a pulp cavity, they are transversed by a number of tubules. The dental formula is I 0/0 C 0/0 Pm 2/2 M 3/3 × 2 for a total of 20 teeth. Anal scent glands are present in both sexes. The senses of sight and hearing are acute, but their eyesight appears to be poor. Adult aardvarks have a head and body length of 100–158 cm (39–62 in.), a tail length of 44–71 cm (17–28 in.), and a shoulder height of 60–65 cm (23–25 in.). Most weigh 50–70 kg (110–153 lb), although some individuals may reach 100 kg (219 lb). Aardvarks exhibit a low basal metabolic rate and a low body temperature of about 34.5◦C (94.1◦F). See DENTITION; SCENT GLAND. The aardvark is one of the world’s great diggers. Its short, thick legs (with four toes on the front feet, five on the rear ones) are armed with long, powerful, sharp-edged, spoon-shaped claws that are intermediate between hoof and nail. The aardvark uses its claws to excavate the extensive burrows in which it dwells. Able to dig with amazing speed, it digs shallow holes when searching for food, larger burrows used for temporary shelter, and extensive tunnel systems in which the young are born. Tunnels may be 13 m (43 ft) long with numerous chambers and several entrances. Only mothers and their young share burrows. Burrows dug by aardvarks may be used by a variety of mammals, including warthogs, hyenas, porcupines, jackals, bat-eared foxes, hares, bats, ground squirrels, and civets, as well as monitor lizards and owls. One bird, the ant-eating chat (Myrmecocicla formicivora), nests in occupied burrows. In agricultural areas, burrows may cause damage to farming equipment and earthen dams. Aardvarks are primarily nocturnal and solitary, and occupy a variety of habitats, including grassy plains, bush country, woodland, and savannah. They appear to prefer sandy soils. The main factor in their distribution is the presence of sufficient quantities of termites and ants. The aardvark is of considerable economic importance in keeping the great hordes of termites in check; if termites are not controlled, they do serious damage. The numbers of termites and ants consumed by the big-bodied aardvark in a single night are staggering. Studies have found that more than 50,000 insects may be consumed by an aardvark in one night. The long, round, thin, sticky, protrusible foot-long tongue is broader than that of the typical anteater or pangolin; together with the well-developed salivary glands, it can retrieve more termites in a shorter period of time than the two others. Since aardvarks do not need to chew their
food, they have no need for incisors or canine teeth. Instead, their stomachs have a muscular pyloric area that serves a function similar to a gizzard in grinding up the food. They occasionally eat beetles, grasshoppers, and the fruit of the wild cucumber, known as the aardvark pumpkin, apparently as a source of water. See ANT; ANTEATER; ISOPTERA. Breeding and development. A single altricial offspring (a young that is born immature and helpless, thus requiring extended parental care) is born after a gestation period of 7–9 months. Births generally occur from July to November. After remaining in the burrow for about 2 weeks, the young aardvark begins accompanying its mother on her nightly foraging excursions. The young can dig for itself when about 6 months old. Sexual maturity is attained at about 2 years of age. Captive aardvarks have lived up to 10 years. Threats. Large flesheaters such as the lion, leopard, and ratel (a fierce, badgerlike animal) frequently prey on aardvarks. Pythons also take the young. African natives relish the aardvark’s meat. They wear its teeth and claws as bracelets or necklaces as good-luck charms and for warding off illnesses. The tough hide may be made into straps and bracelets. Aardvarks have been greatly reduced in numbers and distribution and are now on Appendix 2 of the Convention on International Trade in Endangered Species (CITES). Termites have caused enormous damage to pasture and cereal crops in areas where aardvarks and other insectivorous animals have been exterminated. See ENDANGERED SPECIES. Donald W. Linzey Bibliography. D. Macdonald (ed.), The Encyclopedia of Mammals, Andromeda Oxford, 2001; R. M. Nowak, Walker’s Mammals of the World, 6th ed., Johns Hopkins University Press, 1999.
Abaca One of the strongest of the hard fibers, commercially known as Manila hemp. Abaca is obtained from the leafstalks of a member of the banana family, Musa textilis. The plant resembles the fruiting banana, but is a bit shorter in stature, bears small inedible fruits, and has leaves that stand more erect than those of the banana and are slightly narrower, more pointed, and 5–7 ft (1.5–2 m) long. Relatives of abaca grow wild throughout Southeast Asia, but the plant was domesticated long ago in the southern Philippines. Experiments have succeeded in growing plants yielding fiber in a few other parts of the world, chiefly Middle America, but commercial production has come almost entirely from the Philippines and Borneo. See ZINGIBERALES. Environment and cultivation. Abaca prefers a warm climate with year-round rainfall, high humidity, and absence of strong winds. Soils must always be moist but the plant does not tolerate waterlogging. Abaca grows best on alluvial soils in the southern Philippines and northern Borneo below 1500 ft (457 m) elevation. The plant is best propagated by rootstalk
3
4
Abacus
1
2
(a)
(b)
Appearance of diseased abaca. (a) Affected by mosaic; note ragged leaves and chlorotic areas on leaves. (b) Bunchy top; plant 1 on left is dead and lower leaves of plant 2 have been killed. (USDA)
suckers. There are about 75 varieties grown in the Philippines, grouped into seven categories, each of which varies slightly in height, length, and quality and yield of fiber. Newly planted stock reaches cutting size in 18 months, and a few stalks may be cut from each “hill” every 4 months thereafter for 12– 15 years, after which new plantings are made. Plantings on good soils yield 90–150 lb (41–68 kg) of fiber per acre per year from 2 tons (1.8 metric tons) of stalks and the fiber yield ranges 2–4%. On large commercial holdings plants are grown in solid stands with the new stock spaced 15–20 ft (4.5–6 m) apart, but on small farms plants are scattered and mixed with bamboo and fruit trees. Properties and uses. The fiber ranges 6–14 ft (1.8– 4 m) in strand length, is lustrous, and varies from white to dull yellow. Pre-Spanish Filipinos made clothing with the fine fiber, which today has been replaced by softer, more pliable textiles. As one of the longest and strongest plant fibers, resistant to fresh and salt water, abaca is favored for marine hawsers and other high-strength ropes. Abaca is also used in sackings, mattings, strong papers, and handicraft art goods. Processing and production. The fibers are carried in the outer sections of the leafstalks, and the fleshy material is stripped away mechanically. Traditionally, leafstalks were drawn by hand over a knife blade, under a wooden wedge, and a worker could clean about 15 lb (6.8 kg) of fiber per day. Small power machines strip about 175 lb (79 kg) per worker-day. Large mechanical decorticators save labor and recover more fiber, but the capital outlay is very large.
Diseases. Abaca is affected by several diseases, of which the chief are bunchy top, mosaic, and wilt (see illus.). Bunchy top is caused by a virus spread by the banana aphid (Pentalonia nigronervosa); leaves become smaller and clustered at the plant top, after which the plant dies. Mosaic is also caused by a virus spread by aphids (chiefly Rhopalosiphum nymphaeae and Aphis gossypii); the leaves yellow and dry out as the plant dies. Abaca wilt is caused by a soil or water-borne fungus, chiefly attacking plant roots. Filipino small farmers are not easily brought into disease-control cooperative measures, and postwar epidemics of bunchy top and mosaic have caused serious decline in acreage and production. See PLANT PATHOLOGY. Elton G. Nelson Bibliography. R. E. Huke et al., Shadows on the Land: An Economic Geography of the Philippines, 1963; W. Manshard, Tropical Agriculture, 1975; B. B. Robinson and F. L. Johnson, Abaca: A Cordage Fiber, Agr. Monogr. 21, 1954; J. E. Spencer, The abaca plant and its fiber, Manila hemp, Econ. Bot., 7:195– 213, 1953; F. L. Wernstedt and J. E. Spencer, The Philippine Island World: A Physical, Cultural and Regional Geography, 1967.
Abacus An early mechanical calculator whose design has evolved through the centuries, with two styles in use today. Both the Chinese and the Japanese styles consist of a frame with a crossbeam. They may be made from many different materials, such as wood
Abalone or brass. Rods or wires carrying sliding beads extend vertically through the crossbeam. The Chinese suan pan has two beads above the beam on each rod and five beads below. Each rod of the Japanese soroban carries one bead above and four below. Similar to the abacus in both construction and use, but much larger, are the counting frames used in elementary schools. Braille versions of the abacus are available for use by those without sight. Operation. In working with whole numbers, the rightmost rod represents the ones position, with each rod to the left representing the tens, hundreds, thousands, and so forth, respectively. The beads below the crossbeam represent one of that rod’s units (that is, a one, a ten, a hundred, and so forth), and those above represent five. Beads are moved from the outer position toward the crossbeam when used to represent a number (Fig. 1). By a series of motions using the thumb and forefinger, numbers may be added and subtracted. Addition and subtraction both work from left to right, in contrast to the usual method, using substantial mental arithmetic. For example, to add 384 + 795, the user begins with the beads positioned to represent 384 (Fig. 1). Then to add 795, one lower bead on the fourth rod from the right is moved up, while simultaneously moving the three beads on the third rod down (700 = 1000 − 300). Next, moving to the tens position, the user adds the 90 by the mental method of “add 100 and subtract 10,” thus moving one lower bead up on the hundreds rod and one lower bead down from the center on the tens rod. Finally, to add the 5, the user simply brings one of the upper beads down to the crosspiece. The result of the addition, 384 + 795 = 1179, then appears (Fig. 2). Subtraction is done in a similar manner. In order to subtract 7 from any position, for example, the user subtracts 10 and adds 3. Both multiplication and division can also be performed with the help of the abacus, but these are much more complicated. Use. The abacus, in contrast to the electronic calculator, is simply an aid to mental computation. A well-developed facility with numbers is required in
hundreds
tens
ones
Fig. 1. Beads positioned to represent 384 on a Chinese abacus.
Fig. 2. Beads positioned to represent 384 + 795 = 1179 on a Chinese abacus.
order to use it effectively. For this reason, it is the calculator of choice for teachers in the Far East. The advent of the electronic calculator, while a boon to the scientific community, may have little impact on shop owners who are comfortable with the abacus. See CALCULATORS. Marjorie K. Gross Bibliography. American Printing House for the Blind, Inc., Abacuses, June 1999; C. B. Boyer (revised by U. C. Merzbach), A History of Mathematics, 3d ed., 1991; J. Dilson, The Abacus, St. Martin’s Press, 1995.
Abalone Common name in California for several species of marine gastropod mollusks in the family Haliotidae; also known as ear shells in many tropical and warmtemperate seas, as awabi in Japan, as paua in New Zealand, and as ormers in the Mediterranean region, Spain, and the English Channel. The haliotid limpets constitute the most successful family of the primitive two-gilled gastropods included in the order Archaeogastropoda. They have a broad muscular sucker foot, an ear-shaped limpet shell (up to 30 cm or 12 in. long in some species) with a row of about seven shell apertures (see illus.), and a pair of featherlike gills (aspidobranch ctenidia). They live on rock surfaces in the low intertidal and sublittoral, and are grazing microherbivores, scraping and ingesting algal films with the relatively unspecialized dentition of a rhipidoglossan radula. Haliotids are among the relatively few living gastropods which are zygobranch (with two aspidobranch ctenidia) and diotocardiac (with two auricles in the heart). However, such paired symmetry of pallial, urinogenital, and cardiac structures is almost certainly characteristic of the earliest (ancestral) stocks of true gastropods. In the living forms, which include the rarer pleurotomariids and fissurellids as well as the haliotids, the symmetry of the organs requires a symmetrical respiratory flow through the branchial and pallial structures. The inhalant water current comes in to the mantle cavity from both sides ventrally, and the exhalant current passes out centrally and dorsally. In all of these living zygobranch gastropods
5
6
ABC lipid transporters
Shell of a typical haliotid limpet, or abalone, showing the characteristic curved series of exhalant apertures.
(and all the more numerous fossil forms) there are shell slits or complex apertures in the shell to ensure that the exhalant current, bearing with it both feces and urinogenital products, does not discharge directly over the head. In the haliotids (see illus.) these form a curved series of six or seven exhalant apertures, maintained during shell growth by the oldest (most apical) aperture being filled from below with shell material while a new aperture is formed initially as an indentation in the margin of the shell. Under the penultimate functioning aperture lies the anus, and thus feces are flushed out of the mantle cavity dorsally and well away from the mouth and sense organs. The sexes are separate, with broadcast spawning (again through the older apertures) and external fertilization. The planktonic larvae are typical veligers. In all regions where they occur, flesh of the massive foot of these haliotids is prized as food, and the flattened ear-shaped shells, with characteristically green-blue iridescent inner surfaces (termed motherof-pearl), have been used over many centuries in decorative arts by Maoris, Japanese, Amerindians (Wakashans), and Byzantine Greeks. Collected almost to extinction in some regions for their combined shell beauty and gourmet value, abalones and other haliotids, where protected, have recovered to become again very abundant. See GASTROPODA; LIMPET; PROSOBRANCHIA. W. D. Russell-Hunter
ABC lipid transporters A family of adenosine triphosphate–binding cassette (ABC) transmembrane proteins that use energy to transport various molecules across extracellular and intracellular membranes (cytoplasmic membranes, endoplasmic reticulum, mitochondria, or peroxisomes). ABC lipid transporters consist of two hydrophobic transmembrane domains and two hydrophilic nucleotide-binding folds. These binding folds contain highly conserved sequence mo-
tifs (Walker A and B) and are separated by a linker sequence, also known as the signature (C) motif. ABC transporters are organized as full transporters or half transporters depending on the number of transmembrane domains and nucleotide-binding folds. Half transporters must form homodimers or heterodimers to be functional. ABC transporters are grouped into seven major subfamilies named ABCA to ABCG. The transmembrane domains of the ABC proteins form a chamber within the membrane and act as a “flippase” that gives inner leaflet substrates access to their binding sites. The substrates transported by the ABC proteins can vary from small ions to large polypeptides and polysaccharides. A great number of the ABC transporters are involved with cellular lipid transport and homeostasis, and the defect or absence of these transporters is often correlated with pathologic phenotypes. The table illustrates the major ABC transporters currently associated with lipid transport or homeostasis. ABCA1. ABCA1 (ABC1; CERP, cholesterol efflux regulatory protein) is a member of the subfamily of ABC full-size transporters, containing two highly conserved ATP-binding cassettes (catalytic regions) and two transmembrane domains. ABCA1 is involved in the formation of high-density lipoprotein (HDL) and the cellular efflux of excess cholesterol and phospholipids. The 149-kb gene is located on the human chromosome 9q22-q31. It comprises 50 exons and codes for a calculated 220-kDa protein that is ubiquitously expressed, with highest expression in placenta, liver, lung, adrenal glands, and testes, and modest expression in small intestine, lung, and adipose. The expression of ABCA1 is highly regulated on a transcriptional as well as posttranslational level by several agents, including cholesterol, fatty acids, liver X receptor agonists, cAMP [cyclic adenosine monophosphate], and cytokines. ABCA1 has been identified as the molecular defect in Tangier disease, a rare genetic disorder characterized by nearly complete HDL deficiency, high plasma triglyceride levels, orange tonsils, and increased susceptibility to atherosclerosis. Studies in ABCA1 knockout and transgenic mouse models as well as in-vitro studies have led to major advances in our understanding of the ABCA1-mediated formation of plasma HDL and its role in reverse cholesterol transport. Such transport is considered an antiatherogenic process, in which excess cholesterol is transported from peripheral tissues to the liver in the form of HDL, where it can be excreted from the body directly through the bile or indirectly after conversion into bile acids. ABCA1 is responsible for the initial lipidation of apoA-I, the major apoprotein found on HDL. These nascent or pre-β HDL particles function as acceptors for further lipidation with cholesterol and phospholipids from peripheral tissues, such as lipidladen foam cells found in atherosclerotic lesions. ABCA1 also interacts with serum amyloid A to form large, dense HDL particles associated with pathological conditions such as acute chronic inflammation.
ABC lipid transporters Major ABC transporters currently associated with lipid transport or homeostasis Transporter ABCA1 ABCA3 ABCA7 ABCG1 ABCG2 ABCG4 ABCG5 ABCG8
Function
Disease
Cholesterol and phospholipid transport Surfactant metabolism Cholesterol and phospholipid transport Cholesterol and phospholipid transport Drug resistance Cholesterol transport Sterol transport Sterol transport
Plasma HDL concentrations are regulated primarily by liver-derived (∼70%) and intestine-derived (∼30%) ABCA1 activity and are inversely correlated with coronary artery disease. Complete ABCA1 deficiency in mice did not alter atherosclerosis, but repopulation of ABCA1-deficient mice with wild-type or ABCA1-overexpressing macrophages led to significantly decreased atherosclerosis. Accordingly, selective inactivation of macrophage ABCA1 in mice results in markedly increased atherosclerosis. However, overexpression of ABCA1 in other tissues, including liver, increased plasma cholesterol concentrations of both the antiatherogenic HDL and the proatherogenic non-HDL, and resulted in both proand antiatherogenic effects in different mouse models. Thus, changes in ABCA1 expression alter atherogenesis in a tissue-specific way. Based on these results, upregulation of ABCA1 for therapeutic use may be less desirable, or it may have to target specific tissues to achieve an overall antiatherogenic effect. ABCA3. The ABCA3 gene maps to chromosome 16p13.3, consists of a 1704-amino-acid polypeptide, and is expressed predominantly in lung. ABCA3 gene mutations cause a fatal deficiency of surfactant in newborns. Pulmonary surfactant is mostly composed of lipids that are densely organized into multilamellar structures. These structures are critical to the transition of the lung from an aqueous to an air environment at birth. In adult lung, surfactants are also critical in preserving its homeostasis. Cellular electron micrographs of ABCA3-deficient human lung cells showed lipid accumulation in the multilamellar vesicles, suggesting that ABCA3 plays an important role in the formation of pulmonary surfactant, probably by transporting phospholipids and cholesterol. ABCA7. ABCA7 is a highly expressed transporter in spleen, thymus, lymphoid cells, bone marrow, brain, and trachea. It is mapped to chromosome 19 and has 46 introns, and it is a full-size transporter. ABCA7 is highly homologous to ABCA1 and induces an apolipoprotein-mediated assembly of cholesterol similar to that induced by ABCA1, releasing cholesterol and phospholipids to HDL in an apoA-I–dependent mechanism. In mice, it was reported to promote the apoA-I–dependent release of phospholipids, but not cholesterol. Another observation in ABCA7–deficient females, but not male mice, reported less visceral fat and lower total serum and HDL cholesterol levels than wild-type, gendermatched littermates. Also, in these same mice, the
Tangler disease Neonatal surfactant deficiency Sjogren’s ¨ syndrome Unknown Unknown Unknown Sitosterolemia Sitosterolemia
lack of ABCA7 did not alter cholesterol and phospholipid efflux. ABCG1. ABCG1 (ABC8, or human white gene) is a member of the subfamily of ABC half-size transporters involved in intracellular lipid transport. In order to create a functional full-size transporter, ABCG1 can homodimerize and most probably can also heterodimerize with other half-size transporters. The human ABCG1 spans ∼98 kb and comprises 23 exons, located on chromosome 21q22.3. ABCG1 is expressed in most tissues in mice and humans, with high baseline expression in mouse macrophages, lung, brain, spleen, and other tissues containing macrophages, while in humans the expression is highest in lung, adrenal glands, spleen, heart, placenta, and liver. Several isoforms and splice variants have been found on the RNA level, yet the different protein levels and their functionality remain to be established. Expression of ABCG1 is strongly induced in many tissues by cholesterol loading, by ligands for the nuclear receptors LXRα and β, and by PPARγ -ligands. ABCG1 is responsible for the control of tissue lipid levels and mediates cholesterol and phospholipid efflux to mature HDL or other phospholipid-containing acceptors, particularly in macrophages under conditions of cholesterol excess. The targeted disruption of ABCG1 in mice on a high-fat/high-cholesterol diet resulted in massive accumulation of neutral lipids and phospholipids in hepatocytes and tissue macrophages; however, no changes in plasma lipids accompanied these findings. In accordance with the regulation of cellular lipid homeostasis, the tissues in mice overexpressing human ABCG1 were protected from diet-induced lipid accumulation. ABCG1 was found to be localized perinuclear as well as at the plasma membrane of cells. One potential mechanism for the ABCG1-facilitated lipid transfer is by redistributing membrane cholesterol to cellsurface domains that are accessible to removal by HDL. The role of ABCG1 in cholesterol export of lipid-laden cells seems to complement the action of ABCA1, a full-size transporter with similar tissue expression pattern, which mediates the initial lipidation in the formation of HDL. The nascent HDL particles are substrates for ABCG1-mediated lipid efflux, making ABCG1 part of the reverse cholesterol transport, a protective mechanism against the critical step of macrophage lipid overloading in atherogenesis. ABCG2. Human ABCG2 is also referred to as placenta-specific ABC transporter, breast cancer
7
8
ABC lipid transporters resistance protein, or mitoxantrone resistance– associated protein. (Mitoxantrone is a chemotherapy drug used to treat certain cancers.) It was initially identified as a messenger RNA expressed in placenta. The human ABCG2 gene is located on chromosome 4q22 and encodes a half transporter. In normal human tissues, ABCG2 has been found to be expressed on the apical membrane of trophoblast cells in placenta, the apical membrane of enterocytes, the bile canalicular membrane of hepatocytes, and the apical membrane of lactiferous ducts in the mammary gland. In addition, ABCG2 has been demonstrated to be expressed in a wide variety of hematopoietic cells. It is expected that ABCG2 homodimerizes to acquire transporter activity. ABCG2 transporter functions as a xenobiotic transporter which may play a major role in multidrug resistance. It serves as a cellular defense mechanism in response to mitoxantrone and anthracycline (a chemotherapeutic agent) exposure. Although the function of ABCG2 has been studied extensively in terms of multidrug resistance, the physiological and/or pharmacological functions of this transporter have not been clarified yet. ABCG2 can also transport several dyes, such as rhodamine 123 and Hoechst 33462. Recently, it has been shown to transport steroids such as cholesterol, estradiol, progesterone, and testosterone. It also transports sulfated conjugates of steroids (esterone 3-sulfate, dehydroepiandrosterone sulfate) and certain chlorophyll metabolites. ABCG4. The ABCG4 gene consists of 14 exons that span 12.6 kb on chromosome 11. ABCG4 gene sequence, protein, regulation, and function are very similar to ABCG1. The ABCG4 protein contains 646 amino acids and shares 72% amino acid sequence identity with ABCG1. However, the tissue distribution of ABCG1 and ABCG4 is not identical: ABCG1 is present in many tissues, including macrophages, lung, brain, liver, and spleen, while ABCG4 is mainly expressed in the brain and spleen. ABCG4 expression in the brain has been suggested to promote cholesterol efflux to HDL-like particles present in cerebrospinal fluids. Low level of ABCG4 has also been found in mouse macrophages, and its expression, similarly to ABCG1, was also upregulated by oxysterols (22-OH cholesterol), retinoids (9-cisretinoic acid), and a (liver X receptor)–specific agonist (T0901713). Because of the different pattern of expression, it seems very unlikely that ABCG1 and ABCG4 act as heterodimers. Instead, evidence that ABCG4 acts as homodimer comes from studies that demonstrate that overexpression of ABCG4 stimulated cholesterol efflux to HDL. Both ABCG4 and ABCG1 play a major role in the HDL-mediated cholesterol efflux in liver X receptor– induced macrophages. In-vitro studies showed that ABCG4 expression alone increased cholesterol efflux to both HDL-2 and HDL-3, and when both ABCG4 and ABCG1 were overexpressed, the increase in cholesterol efflux was even higher. ABCG4-mediated cholesterol efflux to LDL and cyclodextrin was also observed, but in lesser degree compared with the
HDL efflux. ABCG4 suppression reduced the cholesterol efflux to HDL, but also decreased ABCG1 messenger RNA expression, probably due to homology in their sequences. These data indicate that ABCG4 alone or in association with ABCG1 may modulate the intracellular concentration of cholesterol by promoting cholesterol efflux to HDL. ABCG5 and ABCG8. ABCG5 and ABCG8 are two half transporters that belong to the ABCG subfamily. They were identified as proteins mutated in sitosterolemia, a rare genetic disorder characterized by increased plasma and tissue levels of plant sterols. Patients with sitosterolemia have an increased absorption of dietary sterols and an impaired ability to secrete cholesterol and plant-derived sterols into bile. Many sitosterolemic individuals deposit cholesterol and plant sterols in the skin, tendons, and coronary arteries and suffer premature coronary artery disease. ABCG5 and ABCG8 genes contain 13 exons and are organized in a head-to-head orientation on chromosome 2 with only 140 base pairs separating their two respective start-transcription sites. In-vivo and in-vitro data indicate that ABCG5 and ABCG8 function as obligate heterodimers to promote sterol excretion into bile. The major regulation of their expression appears to be at the transcriptional level: cholesterol feeding activates the liver X receptor, which increases both ABCG5 and ABCG8 messenger RNA expression. ABCG5/G8 are expressed mainly at the apical surface of hepatocytes and enterocytes, where they limit the accumulation of sterols by promoting sterol excretion into bile and reducing dietary sterol uptake. Overexpression of human ABCG5 and ABCG8 in wild-type mice led to increased cholesterol secretion into bile, and the knockouts of either ABCG5/G8 or ABCG8 alone prevented sterol secretion into bile. However, no changes in liver or plasma cholesterol levels were found. No protection from atherosclerosis was observed in these studies. However, in LDL receptor–deficient mice, which are prone to develop atherosclerosis on a high-fat/high-cholesterol diet, ABCG5/G8 overexpression decreased the absorption of dietary cholesterol, increased the excretion of cholesterol into bile and stool, and decreased cholesterol levels in plasma, leading to a significant reduction in the development of atherosclerotic lesions. See ADENOSINE TRIPHOSPHATE (ATP); ATHEROSCLEROSIS; CHOLESTEROL; GENE; LIPID METABOLISM; LIPOPROTEIN; LIVER; PROTEIN. E. M. Wagner; F. Basso; C. S. Kim; M. J. A. Amar Bibliography. H. B. Brewer, Jr., et al., Regulation of plasma high-density lipoprotein levels by the ABCA1 transporter and the emerging role of high-density lipoprotein in the treatment of cardiovascular disease, Arterioscler. Thromb. Vasc. Biol., 24:1755–1760, 2004; M. Dean, The Human ATPBinding Cassette (ABC) Transporter Superfamily, http://www.ncbi.nlm.nih.gov/books/bv.fcgi?call=& rid=mono 001.TOC&depth=2; T. Janvilisri et al., Sterol transport by the human breast cancer resistance protein (ABCG2) expressed in Lactococcus
Aberration (optics) lactis, J. Biol. Chem., 278:20645–20651, 2003; W. Jessup et al., Roles of ATP binding cassette transporters A1 and G1, scavenger receptor BI and membrane lipid domains in cholesterol export from macrophages. Curr. Opin. Lipidol., 17:247–257, 2006; J. Santamarina-Fojo et al., Complete genomic sequence of the human ABCA1 gene: Analysis of the human and mouse ATP-binding cassette A promoter, Proc. Nat. Acad. Sci. USA, 97:7987–7992, 2000; K. Takahashi et al., ABC proteins: Key molecules for lipid homeostasis, Med. Mol. Morphol., 38:2–12, 2005; N. Wang et al., ATP-binding cassette transporters G1 and G4 mediate cellular cholesterol efflux to high-density lipoproteins, Proc. Nat. Acad. Sci. USA, 101:9774–9779, 2004; K. R. Wilund et al., High-level expression of ABCG5 and ABCG8 attenuates diet-induced hypercholesterolemia and atherosclerosis in Ldlr −/− mice, J. Lipid Res., 45:1429–1436, 2004.
Abdomen A major body division of the vertebrate trunk lying posterior to the thorax; and in mammals, bounded anteriorly by the diaphragm and extending to the pelvis. The diaphragm, found only in mammals, separates the abdominal or peritoneal cavity from the pleural and pericardial cavities of the thorax. In all pulmonate vertebrates (possessing lungs or lunglike organs) other than mammals, the lungs lie in the same cavity with the abdominal viscera, and this cavity is known as the pleuroperitoneal cavity. The large coelomic cavity that occupies the abdomen contains the viscera, such as the stomach, liver, gallbladder, spleen, pancreas, and intestinal tract with their associated nerves and blood vessels. The viscera are contained in the peritoneal sac. This structure is adherent to the abdominal wall as the parietal peritoneum and encompasses all parts of the viscera as an outermost coat, the serosa or visceral peritoneum. Connecting sheets of peritoneum from the body wall to the various organs form the mesenteries, which are always double-walled; blood vessels, lymphatics, nerves, and lymph nodes are located between the two walls of the mesenteries. Other folds of the peritoneum form the omenta. At the lesser curvature of the stomach the serosa continues upward to the liver as the lesser omentum, and from the greater curvature it hangs down as a sheet in front of the intestine, forming the greater omentum. The term abdomen is also applied to a similar major body division of arthropods and other animals. See ANATOMY, REGIONAL. Walter Bock
Aberration (optics) A departure of an optical image-forming system from ideal behavior. Ideally, such a system will produce a unique image point corresponding to each object point. In addition, every straight line in the object
space will have as its corresponding image a unique straight line. A similar one-to-one correspondence will exist between planes in the two spaces. This type of mapping of object space into image space is called a collinear transformation. A paraxial ray trace is used to determine the parameters of the transformation, as well as to locate the ideal image points in the ideal image plane for the system. See GEOMETRICAL OPTICS; OPTICAL IMAGE. When the conditions for a collinear transformation are not met, the departures from that ideal behavior are termed aberrations. They are classified into two general types, monochromatic aberrations and chromatic aberrations. The monochromatic aberrations apply to a single color, or wavelength, of light. The chromatic aberrations are simply the chromatic variation, or variation with wavelength, of the monochromatic aberrations. See CHROMATIC ABERRATION. Aberration measures. The monochromatic aberrations can be described in several ways. Wave aberrations are departures of the geometrical wavefront from a reference sphere with its vertex at the center of the exit pupil and its center of curvature located at the ideal image point. The wave aberration is measured along the ray and is a function of the field height and the pupil coordinates of the reference sphere (Fig. 1). Transverse ray aberrations are measured by the transverse displacement from the ideal image point to the ray intersection with the ideal image plane. Some aberrations are also measured by the longitudinal aberration, which is the displacement along the chief ray of the ray intersection with it. The use
reference sphere (ideal wavefront)
radius of reference sphere (ideal ray)
actual ray wave aberration
actual wavefront
radius of reference sphere (ideal ray)
ideal image point transverse ray aberration
actual ray image plane optical axis
exit pupil Fig. 1. Diagram of the image space of an optical system, showing aberration measures: the wave aberration and the transverse ray aberration.
9
10
Aberration (optics) of this measure will become clear in the discussion of aberration types. Caustics. Another aberrational feature which occurs when the wavefront in the exit pupil is not spherical is the development of a caustic. For a spherical wavefront, the curvature of the wavefront is constant everywhere, and the center of curvature along each ray is at the center of curvature of the wavefront. If the wavefront is not spherical, the curvature is not constant everywhere, and at each point on the wavefront the curvature will be a function of orientation on the surface as well as position of the point. As a function of orientation, the curvature will fluctuate between two extreme values. These two extreme values are called the principal curvatures of the wavefront at that point. The principal centers of curvature lie on the ray, which is normal to the wavefront, and will be at different locations on the ray. The caustic refers to the surfaces which contain the principal centers of curvature for the entire wavefront. It consists of two sheets, one for each of the two principal centers of curvature. For a given type of aberration, one or both sheets may degenerate to a line segment. Otherwise they will be surfaces. The feature which is of greatest interest to the user of the optical system is usually the appearance of the image. If a relatively large number of rays from the same object point and with uniform distribution over the pupil are traced through an optical system, a plot of their intersections with the image plane represents the geometrical image. Such a plot is called a spot diagram, and a set of these is often plotted in planes through focus as well as across the field. Orders of Aberrations The monochromatic aberrations can be decomposed into a series of aberration terms which are ordered according to the degree of dependence they have on the variables, namely, the field height and the pupil coordinates. Each order is determined by the sum of the powers to which the variables are raised in describing the aberration terms. Because of axial symmetry, alternate orders of aberrations are missing from the set. For wave aberrations, odd orders are missing, whereas for transverse ray aberrations, even orders are missing. It is customary in the United States to specify the orders according to the transverse ray designation. Monochromatically, first-order aberrations are identically zero, because they would represent errors in image plane location and magnification, and if the first-order calculation has been properly carried out, these errors do not exist. In any case, they do not destroy the collinear relationship between the object and the image. The lowest order of significance in describing monochromatic aberrations is the third order. Chromatic variations of the first-order properties of the system, however, are not identically zero. They can be determined by first-order ray tracing for the different colors (wavelengths), where the refractive
green blue
red
(a) red green blue
(b) Fig. 2. Chromatic aberration. (a) Longitudinal chromatic aberration. (b) Transverse chromatic aberration.
indices of the media change with color. For a given object plane, the different colors may have conjugate image planes which are separated axially. This is called longitudinal chromatic aberration (Fig. 2a). Moreover, for a given point off axis, the conjugate image points may be separated transversely for the different colors. This is called transverse chromatic aberration (Fig. 2b), or sometimes chromatic difference of magnification. These first-order chromatic aberrations are usually associated with the third-order monochromatic aberrations because they are each the lowest order of aberration of their type requiring correction. The third-order monochromatic aberrations can be divided into two types, those in which the image of a point source remains a point image but the location of the image is in error, and those in which the point image itself is aberrated. Both can coexist, of course. Aberrations of geometry. The first type, the aberrations of geometry, consist of field curvature and distortion. Field curvature. Field curvature is an aberration in which there is a focal shift which varies as a quadratic function of field height, resulting in the in-focus images lying on a curved surface. If this aberration alone were present, the images would be of good quality on this curved surface, but the collinear condition of plane-to-plane correspondence would not be satisfied. Distortion. Distortion, on the other hand, is an aberration in which the images lie in a plane, but they are displaced radially from their ideal positions in the image plane, and this displacement is a cubic function of the field height. This means that any straight
Aberration (optics)
(a)
(b)
Fig. 3. Distortion. (a) Barrel distortion. (b) Pincushion distortion.
line in the object plane not passing through the center of the field will have an image which is a curved line, thereby violating the condition of straight-line to straight-line correspondence. For example, if the object is a square centered in the field, the points at the corners of the image are disproportionally displaced from their ideal positions in comparison with the midpoints of the sides. If the displacements are toward the center of the field, the sides of the figure are convex; this is called barrel distortion (Fig. 3a). If the displacements are away from the center, the sides of the figure are concave; this is called pincushion distortion (Fig. 3b). Aberrations of point images. There are three thirdorder aberrations in which the point images themselves are aberrated: spherical aberration, coma, and astigmatism. Spherical aberration. Spherical aberration is constant over the field. It is the only monochromatic aberration which does not vanish on axis, and it is the axial case which is easiest to understand. The wave aberration function (Fig. 4a) is a figure of revolution which varies as the fourth power of the radius in the pupil. The wavefront itself has this wave aberration function added to the reference sphere centered on the ideal image. The rays (Fig. 4b) from any circular zone in the wavefront come to a common zonal focus on the axis, but the position of this focus shifts axially as the zone radius increases. This zonal focal shift increases quadratically with the zone radius to a maximum from the rays from the marginal zone. This axial shift is longitudinal spherical aberration. The magnitude of the spherical aberration can be measured by the distance from the paraxial focus to the marginal focus. The principal curvatures of any point on the wavefront are oriented tangential and perpendicular to the zone containing the point. The curvatures which are oriented tangential to the zone have their centers on the axis, so the caustic sheet for these consists of the straight line segment extending from the paraxial focus, and is degenerate. The other set of principal centers of curvature lie on a trumpet-shaped surface concentric with the axis (Fig. 5a). This second sheet is the envelope of the rays. They are all tangent to it, and the point of tangency for each ray is at the second principal center of curvature for the ray (Fig. 5b). It is clear that the image formed in the presence of spherical aberration (Fig. 4c) does not have a well-
defined focus, although the concentration of light outside the caustic region is everywhere worse than it is inside. Moreover, the light distribution in the image is asymmetric with respect to focal position, so the precise selection of the best focus depends on the criterion chosen. The smallest circle which can contain all the rays occurs one-quarter of the distance from the marginal focus to the paraxial focus, the image which has the smallest second moment lies one-third of the way from the marginal focus to the paraxial focus, and the image for which the variance of the wave aberration function is a minimum lies halfway between the marginal and paraxial foci. Coma. Coma is an aberration which varies as a linear function of field height. It can exist only for offaxis field heights, and as is true for all the aberrations, it is symmetrical with respect to the meridional plane containing the ideal image point. Each zone in its wave aberration function (Fig. 6a) is a circle, but each circle is tilted about an axis perpendicular to the meridional plane, the magnitude of the tilt increasing with the cube of the radius of the zone. The chief ray (Fig. 6b) passes through the ideal
(a)
zone in pupil
image plane
zonal focus
(b)
ideal image point (paraxial focus)
(c) Fig. 4. System with spherical aberration. (a) Wave aberration function. (b) Rays. (c) Spot diagrams through foci showing transverse ray aberration patterns for a square grid of rays in the exit pupil.
11
12
Aberration (optics) end of caustic
marginal rays
(a)
paraxial marginal focus focus minimum circle
external caustic
axial caustic
(b) Fig. 5. Caustics in a system with spherical aberration. (a) Diagram of caustics and other major features. (b) Rays whose envelope forms an external caustic.
image point, but the rays from any zone intersect the image plane in a circle, the center of which is displaced from the ideal image point by an amount equal to the diameter of the circle. The diameter increases as the cube of the radius of the corresponding zone in the pupil. The circles for the various zones are all tangent to two straight lines intersecting at the ideal image point and making a 60◦ angle with each other. The resulting figure (Fig. 6c) resembles an arrowhead which points toward or away from the center of the field, depending on the sign of the aberration. The upper and lower marginal rays in the meridional plane intersect each other at one point in the circle for the marginal zone. This point is the one most distant from the chief ray intersection with the image plane. The transverse distance between these points is a measure of the magnitude of the coma. Astigmatism. Astigmatism is an aberration which varies as the square of the field height. The wave aberration function (Fig. 7a) is a quadratic cylinder which varies only in the direction of the meridional plane and is constant in the direction perpendicular to the meridional plane. When the wave aberration function is added to the reference sphere, it is clear that the principal curvatures for any point in the wavefront are oriented perpendicular and parallel to the meridional plane, and moreover, although they are different from each other, each type is constant over the wavefront. Therefore, the caustic sheets both degenerate to lines perpendicular to and in the meridional plane in the image region, but the two lines are separated along the chief ray. All of the rays must pass through both of these lines, so they are identified as the astigmatic foci. The
astigmatic focus which is in the meridional plane is called the sagittal focus, and the one perpendicular to the meridional plane is called the tangential focus. For a given zone in the pupil, all of the rays (Fig. 7b) will of course pass through the two astigmatic foci, but in between they will intersect an image plane in an ellipse, and halfway between the foci they will describe a circle (Fig. 7c). Thus, only halfway between the two astigmatic foci will the image be isotropic. It is also here that the second moment is a minimum, and the wave aberration variance is a minimum as well. This image is called the medial image. Since astigmatism varies as the square of the field height, the separation of the foci varies as the square of the field height as well. Thus, even if one set of foci, say the sagittal, lies in a plane, the medial and tangential foci will lie on curved surfaces. If the field curvature is also present, all three lie on curved surfaces. The longitudinal distance along the chief ray from the sagittal focus to the tangential focus is a measure of the astigmatism. The above description of the third-order aberrations applies to each in the absence of the other aberrations. In general, more than one aberration will be present, so that the situation is more complicated. The types of symmetry appropriate to each aberration will disclose its presence in the image.
zone in pupil (a)
ideal image point
image plane
(b)
(c) Fig. 6. System with coma. (a) Wave aberration function. (b) Rays. (c) Spot diagrams through foci, showing transverse ray aberration patterns for a square grid of rays in the exit pupil.
Aberration (optics) used, although in principle they are available. In fact, many optical designers use the third order as a guide in developing the early stages of a design, and then go directly to real ray tracing, using the transverse ray aberrations of real rays without decomposition to orders. However, the insights gained by using the fifth-order aberrations can be very useful. Origin of Aberrations
(a) zone in pupil tangential image
ideal image point sagittal image image plane
(b)
tangential focus
medial focus
sagittal focus
(c) Fig. 7. System with astigmatism. (a) Wave aberration function. (b) Rays. (c) Spot diagrams through foci, showing transverse aberration patterns for a square grid of rays in the exit pupil.
Higher-order aberrations. The next order of aberration for the chromatic aberrations consists of the chromatic variation of the third-order aberrations. Some of these have been given their own names; for example, the chromatic variation of spherical aberration is called spherochromatism. Monochromatic aberrations of the next order are called fifth-order aberrations. Most of the terms are similar to the third-order aberrations, but with a higher power dependence on the field or on the aperture. Field curvature, distortion, and astigmatism have a higher power dependence on the field, whereas spherical aberration and coma have a higher power dependence on the aperture. In addition, there are two new aberration types, called oblique spherical aberration and elliptical coma. These are not directly related to the third-order terms. Expansions beyond the fifth order are seldom
Each surface in an optical system introduces aberrations as the beam passes through the system. The aberrations of the entire system consist of the sum of the surface contributions, some of which may be positive and others negative. The challenge of optical design is to balance these contributions so that the total aberrations of the system are tolerably small. In a well-corrected system the individual surface contributions are many times larger than the tolerance value, so that the balance is rather delicate, and the optical system must be made with a high degree of precision. Insight as to where the aberrations come from can be gained by considering how the aberrations are generated at a single refracting surface. Although the center of curvature of a spherical surface lies on the optical axis of the system, it does not in itself have an axis. If, for the moment, the spherical surface is assumed to be complete, and the fact that the entrance pupil for the surface will limit the beam incident on it is ignored, then for every object point there is a local axis which is the line connecting the object point with the center of curvature. All possible rays from the object point which can be refracted by the surface will be symmetrically disposed about this local axis, and the image will in general suffer from spherical aberration referred to this local axis. A small pencil of rays about this axis (locally paraxial) will form a first-order image according to the rules of paraxial ray tracing. If the first-order imagery of all the points lying in the object plane is treated in this manner, it is found that the surface containing the images is a curved surface. The ideal image surface is a plane passing through the image on the optical axis of the system, and thus the refracting surface introduces field curvature. The curvature of this field is called the Petzval curvature. In addition to this monochromatic behavior of the refracting surface, the variation in the index of refraction with color will introduce changes in both the first-order imagery and the spherical aberration. This is where the chromatic aberrations come from. Thus there are fundamentally only three processes operating in the creation of aberrations by a spherical refracting surface. These result in spherical aberration, field curvature, and longitudinal chromatic aberration referred to the local axis of each image. In fact, if the entrance pupil for the surface is at the center of curvature, these will be the only aberrations that are contributed to the system by this refracting surface. In general, the pupil will not be located at the center of curvature of the surface. For any off-axis object point, the ray which passes through the center of the
13
14
Aberration of light pupil will be the chief ray, and it will not coincide with the local axis. The size of the pupil will limit the beam which can actually pass through the surface, and the beam will therefore be an eccentric portion of the otherwise spherically aberrated beam. Since an aberration expansion decomposes the wave aberration function about an origin located by the chief ray, the eccentric and asymmetric portion of an otherwise purely spherically aberrated wave gives rise to the field-dependent aberrations, because the eccentricity is proportional to the field height of the object point. In this manner the aberrations which arise at each surface in the optical system, and therefore the total aberrations of the system, can be accounted for. See LENS (OPTICS); OPTICAL SURFACES. Roland V. Shack Bibliography. M. Born and E. Wolf, Principles of Optics, 7th ed., 1999; B. D. Guenther, Modern Optics, 1990; R. Kingslake, Optical System Design, 1983; K. Narigi and Y. Matsui, Fundamentals of Practical Aberration Theory, 1993; W. Smith, Modern Optical Engineering, 2d ed., 1990; W. T. Welford, Aberrations of Optical Systems, 1986.
Aberration of light The apparent change in direction of a source of light caused by an observer’s component of motion perpendicular to the impinging rays. To visualize the effect, first imagine a stationary telescope (illus. a) aimed at a luminous source such as a star, with photons traveling concentrically down the tube to an image at the center of the focal plane. Next give the telescope a component of motion perpendicular to the incoming rays (illus. b). Photons
( true
)
(apparent)
direction of star α photons entering this aperture
Abrasive
reach this focal surface
(a)
(b)
passing the objective require a finite time to travel the length of the tube. During this time the telescope has moved a short distance, causing the photons to reach a spot on the focal plane displaced from the former image position. To return the image to the center, the telescope must be tilted in the direction of motion by an amount sufficient to ensure that the photons once again come concentrically down the tube in its frame of reference (illus. c). The necessary tilt angle α is given by tan α = v/c, where v is the component of velocity perpendicular to the incoming light and c is the velocity of light. (An analogy illustrating aberration is the experience that, in order for the feet to remain dry while walking through vertically falling rain, it is necessary to tilt an umbrella substantially forward.) Aberration of light was discovered by the English astronomer James Bradley in 1725 while searching unsuccessfully for parallax of nearby stars, using the Earth’s orbital diameter as baseline. He found instead the necessity to compensate continuously for the Earth’s velocity in its elliptical orbit, an effect about a hundred times greater than the parallax of typical nearby stars. This discovery provided the first direct physical confirmation of the Copernican theory. Annual aberration causes stars to appear to describe ellipses whose major axes are the same length, the semimajor axis (usually expressed in arc-seconds) being the constant of aberration α = 20.49552. The minor axis of the ellipse depends on the star’s ecliptic latitude β, and is given by 2α sin β. The small diurnal aberration arising from the Earth’s axial rotation depends on geographic latitude φ, and is given by 0.31 cos φ. A second important application of aberration has been its clear-cut demonstration that, as is axiomatic to special relativity, light reaching the Earth has a velocity unaffected by the relative motion of the source toward or away from the Earth. This is shown by the fact that the aberration effect is the same for all celestial objects, including some quasars with apparent recessional velocities approaching the speed of light. See EARTH ROTATION AND ORBITAL MOTION; LIGHT; PARALLAX (ASTRONOMY); QUASAR; RELATIVITY. Harlan J. Smith
(c)
Demonstration of aberration. (a) Fixed telescope; photons form image at center of focal plane. (b) Moving telescope; image is displaced from center. (c) Tilted moving telescope is required to restore image to center.
A material of extreme hardness that is used to shape other materials by a grinding or abrading action. Abrasive materials may be used either as loose grains, as grinding wheels, or as coatings on cloth or paper. They may be formed into ceramic cutting tools that are used for machining metal in the same way that ordinary machine tools are used. Because of their superior hardness and refractory properties, they have advantages in speed of operation, depth of cut, and smoothness of finish. Abrasive products are used for cleaning and machining all types of metal, for grinding and polishing ceramics and glass, for grinding logs to paper pulp, for cutting metals, glass, and cement, and for
Abscisic acid manufacturing many miscellaneous products such as brake linings and nonslip floor tile. Abrasive materials. These may be classified in two groups, the natural and the synthetic (manufactured). The latter are by far the more extensively used, but in some specific applications natural materials still dominate. The important natural abrasives are diamond (the hardest known material), corundum (a relatively pure, natural aluminum oxide, Al2O3), and emery (a less pure Al2O3 with considerable amounts of iron). The last of these has been replaced to a great extent by synthetic materials. Other natural abrasives are garnet, an aluminosilicate mineral; feldspar, used in household cleansers; calcined clay; lime; chalk; and silica, SiO2, in its many forms—sandstone, sand (for grinding plate glass), flint, and diatomite. The synthetic abrasive materials are silicon carbide, SiC; aluminum oxide, Al2O3; titanium carbide, TiC; and boron carbide, B4C. The synthesis of diamond puts this material in the category of manufactured abrasives. There are other carbides, as well as nitrides and cermets, which can be classified as abrasives but their use is special and limited. The properties of hardness and high melting point are related, and many abrasive materials are also good refractories. See DIAMOND; REFRACTORY. Silicon carbide is still made in much the same way that E. G. Acheson first made it in 1891. The principal ingredients are pure sand, SiO2; coke (carbon); sawdust (to burn and provide vent holes for the escape of the gaseous products); and salt (to react with the impurities and form volatile compounds). The batch is placed in a troughlike furnace, up to 60 ft (18 m) long, with a graphite electrode at each end. The charge is heated by an electric current, a core of carbon or graphite being used to carry the current. Temperatures up to 4400◦F (2400◦C) are reached. The net reaction is SiO2 + 3C → SiC + 2CO, but the details of the reaction are complicated and not thoroughly understood. After the furnace has cooled, the sides are removed and the usable silicon carbide picked out, crushed, and sized. Abrasive aluminum oxide is made from calcined bauxite (a mixture of aluminum hydrates) or purified alumina by fusion in an electric-arc furnace with a water-cooled metal shell; the solid alumina that is around the edge of the furnace acts as the refractory. Various grades of each type of synthetic abrasive are distinguishable by properties such as color, toughness, and friability. These differences are caused by variation in purity of materials and methods of processing. The sized abrasive may be used as loose grains, as coatings on paper or cloth to make sandpaper and emery cloth, or as grains for bonding into wheels. Abrasive wheels. A variety of bonds are used in making abrasive wheels: vitrified or ceramic, essentially a glass or glass plus crystals; sodium silicate; rubber; resinoid; shellac; and oxychloride. Each type of bond has its advantages. The more rigid ceramic bond is better for precision-grinding operations, and the tougher, resilient bonds, such as resinoid or rub-
ber, are better for snagging and cutting operations. Ceramic-bonded wheels are made by mixing the graded abrasive and binder, pressing to general size and shape, firing, and truing or finishing by grinding to exact dimensions. A high-speed rotation test is given to check the soundness of the wheel. Grinding wheels are specified by abrasive type, grain size (grit), grade or hardness, and bond type; in addition, the structure or porosity may be indicated. The term hardness as applied to a wheel refers to its behavior in use and not to the hardness of the abrasive material itself. The wheel is a three-component system of abrasive, bond, and air; the hardness is a complex function of the type and amount of bond and of the density of the wheel. Literally thousands of types of wheels are made with different combinations of characteristics, not to mention the multitude of sizes and shapes available; therefore, selecting the best grinding wheel for a given job is not simple. When the abrasive grains of a wheel become slightly dulled by use, the stresses in the grinding operation should increase enough to tear the grain from the wheel to expose a new cutting grain. Thus, too soft a wheel wears too fast, losing grains before they are dulled, whereas too hard a wheel develops a smooth, glazed surface that will not cut. In either case, efficiency drops. Hardness tests. The hardness of materials is roughly indicated by the Mohs scale, in which 10 minerals were selected, arranged in order of hardness, and numbered from 1 to 10 (softest to hardest); diamond has a hardness number of 10 and talc a hardness number of 1. However, since this scale is not quantitative (for example, aluminum oxide, 9 on the Mohs scale, is not 10% softer than diamond), and since most abrasive materials fall in the region at the top of the scale, other hardness scales have been developed. One is the Knoop indentation test, in which a diamond pyramid of specified shape is pressed under a definite load into the material to be tested; the size of the indentation is taken as a measure of the material’s hardness. See HARDNESS SCALES. Ceramic cutting tools are generally made of aluminum oxide by hot-pressing and grinding to final size. See CERAMICS. John F. McMahon Bibliography. S. Krar and E. Ratterman, Superabrasives: Grinding and Machining with Carbon and Diamond,1990; J. Wang, Abrasive Technology: Current Development and Applications I,1999.
Abscisic acid One of the five major plant hormones that play an important role in plant growth and development. Abscisic acid (ABA) is ubiquitous in higher plants; the highest levels are found in young leaves and developing fruits and seeds. Abscisic acid is also produced by certain algae and by several phytopathogenic fungi. It has been found in the brains of several mammals; however, animals probably do not produce abscisic acid but acquire it from plants in the diet.
15
16
Abscisic acid Chemistry and biosynthesis. The chemical structure of (+)-abscisic acid is shown here; naturally 8'
6
9' 5
5'
6'
2
OH
4'
O
3'
3
1'
2'
7'
4
COOH
occurring abscisic acid is dextrorotary, hence the (+). It contains 15 carbon atoms and is a weak organic acid, as are two other plant hormones, auxin and gibberellin. See AUXIN; GIBBERELLIN. Abscisic acid is a terpenoid. All chemicals belonging to this class of compounds have mevalonic acid as a common precursor and are built up of five carbon units (isoprenoid units). Other examples of terpenoids in plants are essential oils, gibberellins, steroids, carotenoids, and natural rubber. See CAROTENOID; ESSENTIAL OILS; RUBBER; STEROID. In fungi, abscisic acid is synthesized by the socalled direct pathway, which means that abscisic acid molecules are built up by combining three isoprenoid units. The biosynthetic pathway in higher plants is not clear, but presumably it is an indirect pathway that involves the breakdown of a larger precursor molecule. Several lines of evidence suggest that this precursor may be a carotenoid: mutants lacking carotenoids also fail to produce abscisic acid, and inhibitors that block carotenoid biosynthesis at the same time inhibit abscisic acid formation. Furthermore, labeling studies with 18O2 indicate that one 18O atom is rapidly incorporated into the carboxyl group of abscisic acid, whereas the ring oxygens become labeled much more slowly. These observations are in accordance with the hypothesis that abscisic acid is derived from a large precursor molecule in which the ring oxygens are already present and an oxygen atom is introduced into the carboxyl group by oxidative cleavage. Physiological roles. Abscisic acid was originally discovered by F. T. Addicott as a factor that accelerates abscission of young cotton fruits (hence the original name abscisin), and by P. F. Wareing as a hormone (called dormin) that causes bud dormancy in woody species. However, subsequent work established that abscisic acid has little or no effect on abscission of leaves and fruits and that it is not an endogenous regulator of bud dormancy in woody species. See ABSCISSION; DORMANCY. Several approaches have been used to determine the roles of abscisic acid in plants. Application of abscisic acid results in a multitude of responses, but most of these may be of a pharmacological nature rather than reflect the roles of endogenous abscisic acid. The study of mutants that are deficient in abscisic acid, such as the flacca mutant of tomato and droopy in potato, has been more informative; both mutants synthesize much less abscisic acid than wildtype plants. These mutants wilt readily because they do not close their stomata. When they are treated with abscisic acid, the normal phenotype is restored,
thus demonstrating that abscisic acid has a key role in preventing excessive water loss. Abscisic acid also functions in growth of roots and shoots, in heterophylly in aquatic plants, and in preventing premature germination of developing embryos. On the other hand, there is now increasing evidence against a role for abscisic acid in gravitropic curvature of roots. Seedlings in which abscisic acid synthesis is inhibited or mutants in which levels of abscisic acid are reduced, show the same gravitropic response of their roots as control seedlings. Stress. Biosynthesis of abscisic acid is greatly accelerated under stress conditions. For example, in a wilted leaf the abscisic acid content can increase as much as 40 times the original level over a 4- to 5-h period. Following rehydration of such a leaf, the abscisic acid level returns to that of a turgid leaf after 5 to 6 h. Thus, the water status of the leaf, in particular the loss of turgor, regulates the rate at which abscisic acid is synthesized and degraded. The excess abscisic acid is either oxidized to phaseic acid, which in turn may be converted to dihydrophaseic acid, or it may be conjugated to glucose to give a complex that is stored in the vacuole. A very rapid response to an increase in abscisic acid in a leaf is the closure of stomata: abscisic acid affects the permeability of the guard cells, which results in loss of solutes and turgor, so that the stomata close and the plant is protected from desiccation. Abscisic acid is called a stress hormone, because it can ameliorate the effects of various stresses (drought, cold, salinity) to which plants are exposed. Freezing tolerance can be induced in certain plants by application of abscisic acid. See COLD HARDINESS (PLANT); PLANT-WATER RELATIONS; PLANTS, SALINE ENVIRONMENTS OF. Growth. Abscisic acid inhibits the growth of stems and expansion of young leaves by decreasing cellwall loosening, an effect opposite to that caused by auxin. On the other hand, growth of roots may be promoted by abscisic acid under certain conditions. These opposite effects of abscisic acid on the growth of roots and shoots are advantageous for the survival of plants under water stress conditions. By inhibiting shoot growth, turgor pressure is maintained, whereas a stimulation of root growth by abscisic acid enlarges the root system, thus increasing water uptake. See PLANT GROWTH. Heterophylly. In aquatic plants two distinct types of leaves are produced on the individual. Submerged leaves are highly dissected or linear without stomata, whereas aerial leaves are entire and have stomata. When abscisic acid is added to the medium, aerial leaves with stomata are produced on submerged shoots. It is conceivable that under natural conditions shoot tips emerging from the water experience water stress, which causes production of abscisic acid, which in turn results in a change in leaf morphology. Embryogenesis, seed maturation, and germination. Following fertilization in higher plants, the ovule develops into the seed with the embryo inside. During
Abscission late embryogenesis the seed desiccates and the embryo becomes dormant. However, if young embryos are excised from the seed prior to the onset of dormancy and cultured in the laboratory, they will germinate. This precocious germination of excised, immature embryos can be prevented by application of abscisic acid. There is convincing evidence that abscisic acid produced in the seed plays an important role in preventing precocious germination early in embryogenesis. Mutants deficient in abscisic acid all show the phenomenon of vivipary, which means that the embryo germinates precociously while the seed is still attached to the mother plant. Elegant genetic experiments with seeds of the small crucifer Arabidopsis thaliana have demonstrated that it is abscisic acid produced by the embryo, not by the maternal tissues of the seed, that controls dormancy of the embryo. Abscisic acid also induces the accumulation of specific proteins during late embryogenesis, but it is not clear at present that such proteins play a role in the onset of embryo dormancy. See PLANT MORPHOGENESIS; SEED. During germination the reserve foods (carbohydrate, protein, and lipid) in the seed are degraded and mobilized to the young seedling. The hormonal regulation of enzymes involved in the mobilization of these reserves has been studied extensively in cereal grains, particularly in barley. Various hydrolytic enzymes, for example, α-amylase, are synthesized by aleurone cells in response to gibberellin produced by the germinating embryo. It has long been known that applied abscisic acid can inhibit the synthesis of these hydrolytic enzymes. Also, it has been found that abscisic acid produced in germinating seeds during water stress can have the same effect. So, the germinating seed not only has the means to turn on the metabolic machinery to mobilize food reserves but can also turn it off during conditions unfavorable for growth. See PLANT PHYSIOLOGY. Mode of action. Two types of responses to abscisic acid can be distinguished. The rapid response, such as stomatal closure, occurs within a few minutes and can be induced only by naturally occurring (+)-abscisic acid. Since abscisic acid is active on guard cells over a wide pH range, while it is taken up by the cells only at low pH, the site of action is likely on the outer surface of the plasma membrane. Stomatal closure is caused by the rapid release of solutes, in particular the potassium ion (K+) and malate, from the guard cells. Thus, in the rapid response abscisic acid affects the permeability of membranes. The slow responses, which take hours or even days to become apparent, involve gene expression and can be induced equally well by (+)-abscisic acid and synthetic (−)-abscisic acid. This implies that the fast and slow responses have different receptor sites. The slow responses involve inhibition of the synthesis of certain enzymes, as well as induction of synthesis of a new group of proteins, the so-called abscisic acid–inducible proteins. These proteins appear during late embryogenesis when the seed begins to desiccate, and during water stress in other
organs of the plant. The genes for these proteins are expressed only in the presence of abscisic acid or after water stress. When examined at the level of messenger ribonucleic acid, the effects of abscisic acid and water stress are not cumulative, indicating that both act by way of the same mechanism. Changes in protein pattern induced by water stress persist but are rapidly reversed by rehydration. The function of the abscisic acid–inducible proteins is not known, but it is thought that they have a function in drought tolerance. See PLANT HORMONES. Jan A. D. Zeevaart Bibliography. F. T. Addicott (ed.), Abscisic Acid, 1983; J. A. D. Zeevaart and R. A. Creelman, Metabolism and physiology of abscisic acid, Annu. Rev. Plant Physiol. Plant Mol. Biol., 39:439–473, 1988.
Abscission The process whereby a plant sheds one of its parts. Leaves, flowers, seeds, and fruits are parts commonly abscised. Almost any plant part, such as very small buds and bracts to branches several inches in diameter, and scales or sheets of bark, may be abscised by some species. However, other species, including many annual plants, may show little abscission, especially of leaves. Abscission may be of value to the plant in several ways. It can be a process of self-pruning, removing injured, diseased, or senescent parts. It permits the dispersal of seeds and other reproductive structures. It facilitates the recycling of mineral nutrients to the soil. It functions to maintain homeostasis in the plant, keeping in balance leaves and roots, and vegetative and reproductive parts. For example, trees living in climates that have a pronounced dry season are often adapted to abscising all or part of their leaves during the dry period, thus keeping the trees in balance with the water of the environment. Flowers and young fruit are usually produced in much larger numbers than the plant is capable of developing into mature fruit. Abscission of flowers and young fruits is the means whereby the plant keeps the number of fruits in balance with the ability of the leaves to nourish them. The selective abscission of branches and branchlets is an important factor in determining the architecture of a number of tree species. Mechanisms. In most plants the process of abscission is restricted to an abscission zone at the base of an organ (see illus.); here separation is brought about by the disintegration of the walls of a special layer of cells, the separation layer. Less than an hour may be required to complete the entire process for some flower petals, but for most plant parts several days are required for completion of abscission. The portion of the abscission zone which remains on the plant commonly develops into a corky protective layer that becomes continuous with the cork of the stem. In general, a healthy organ inhibits its own abscission. Only after it is injured, diseased, or senescent
17
18
Abscission
axillary bud separation layer protective layer
(a)
abscission zone (b)
cork (c)
Diagrams of the abscission zone of a leaf. (a) A leaf with the abscission zone indicated at the base of the petiole. (b) The abscission zone layers shortly before abscission and (c) the layers after abscission.
is an organ abscised. In many instances abscission is also a correlation phenomenon, resulting from events elsewhere in the plant. For example, the abscission of flower petals usually follows closely after pollination and fertilization; and in some broadleaved evergreen trees, such as live oaks, the camphor tree, and Magnolia grandiflora, the major flush of leaf abscission follows closely after the appearance of new buds and leaves in the spring. Abscission is a typical physiological process requiring favorable temperatures, oxygen, and energyyielding respiration. The energy is required for the synthesis of hydrolytic enzymes. These enzymes function to soften the pectins and other constituents of the cell wall. When the cell walls of the separation layer are sufficiently weakened, the leaf (or other organ) falls from the plant by its own weight. Abscission takes place in the lower plants (mosses, ferns, algae, and fungi) in the same way as in the higher plants. Enzymes soften and weaken the pectinlike materials that cement cells together. The cementing substances of higher and lower plants are quite similar in that all are composed of branching polysaccharides. Abscission in the lower plants is less conspicuous because of their simpler structure. However, propagules, such as buds and plantlets, are regularly abscised. In the algae, fragmentation by abscission of segments of filaments is a common method of vegetative propagation. Yeast buds, most mold spores, and bacterial colonies also are separated by the dissolution of polysaccharide cementing substances. Evolution. The fossil record shows that abscission has been taking place for a very long time. The record contains many instances of the abscission of fruits, seeds, and seedlike structures going back millions of years. Leaf abscission became evident soon after leaves evolved in the Devonian Period, some 400 million years ago. Fossils of the oldest known green plants, blue-green algae of the Precambrian Period, 850 million years ago, show that their filaments were abscising (fragmenting), as do present-day bluegreen algae. From the evidence summarized above, it is reasonable to conclude that the evolution of abscission began when two cells first separated. See FOSSIL SEEDS AND FRUITS; PALEOBOTANY.
Factors affecting abscission. Observations of higher plants show that abscission is affected by a number of environmental and internal factors. It can be initiated or accelerated by extremes of temperature, such as frost; by extremes of moisture, such as drought or flooding; by deficiency of mineral elements in the soil, particularly nitrogen, calcium, magnesium, potassium, and zinc; by the shortening photoperiod in the fall; by oxygen concentrations above 20%; by volatile chemicals, such as ethylene and air pollutants; by moderately toxic chemicals; and by the feeding of some injurious insects. Abscission can be retarded or inhibited by excessive nitrogen from the soil; by high carbohydrate levels in the plant; by oxygen concentrations below 20%; and by application of growth-regulator chemicals, such as naphthaleneacetic acid. Hormones. Although a number of internal chemical changes regularly precede the abscission of an organ, at most only a few of the changes appear to be directly related to the process of abscission. Of the changes, a decrease in the level of the growth hormone auxin appears to be the most important. Auxin applied experimentally to the distal (organ) side of an abscission zone retards abscission, while auxin applied to the proximal (stem) side accelerates abscission. From this and related evidence, it is considered that the gradient of auxin across the abscission zone is the major internal factor controlling abscission. Further, many other factors affecting abscission appear to act through effects on the auxin gradient. See AUXIN. The gibberellins are growth hormones which influence a number of plant processes, including abscission. When applied to young fruits or to leaves, they tend to promote growth, delay maturation, and thereby indirectly prevent or delay abscission. See GIBBERELLIN. Abscisic acid is a different type of hormone and has the ability to promote abscission and senescence and to retard growth. Like auxin, it is synthesized in fruits and leaves; however, it tends to counteract the effects of the growth-promoting hormones, auxin, gibberellin, and cytokinin. See ABSCISIC ACID; CYTOKININS. Another hormonal chemical is the gas ethylene. Small amounts of ethylene have profound effects on the growth of plants and can distort and reduce growth and promote senescence and abscission. See ETHYLENE. Thus the process of abscission is strongly influenced by interaction of several plant hormones. Frequently the hormones that tend to retard or delay abscission (auxin, gibberellin) are counteracted by the hormones that tend to promote and accelerate abscission (abscisic acid, ethylene). These interactions are further influenced by many environmental and external factors, as mentioned above. Depending on circumstances, any one of several factors may at times be the key factor controlling the entire process of abscission. See PLANT HORMONES. Agricultural applications. In agricultural practice, abscission is delayed or prevented by keeping soil
Absolute zero moisture and soil nitrogen at high levels. High soil moisture helps to keep abscisic acid at relatively low levels, and high soil nitrogen leads to high levels of auxin in the plant. Auxin-type chemical regulators are frequently used to retard abscission. Acceleration of abscission can be obtained by low soil nitrogen (leading to low auxin), by moisture stress (leading to high abscisic acid), by ethylene-releasing chemicals, or by mildly toxic chemicals. To delay abscission of leaves, fruits, and other plant parts, the auxin-type regulator, naphthaleneacetic acid, or related compounds are commonly used. Some notable applications are the spraying of apple orchards and pear orchards to prevent abscission of nearly mature fruit and the dipping of cut holly to prevent abscission of berries and leaves during shipment. To accelerate abscission a variety of chemicals find use in one or another practice. Cotton is chemically defoliated to facilitate mechanical harvest. To promote leaf abscission, chemicals such as magnesium chlorate, tributylphosphorotrithioate, and merphos have been widely used. In some regions, young nursery plants are defoliated to facilitate earlier digging and shipping. Certain varieties of apples, peaches, and similar fruits must be thinned to ensure that the remaining fruit can reach marketable size. Naphthaleneacetamide and related compounds are the most frequently used to thin fruit. Although these chemicals are of the auxin type, the dosages applied to the flowers and young fruit are somewhat toxic and lead to abortion and abscission of a portion of the young fruit. The slightly toxic insecticide carbaryl is also used in fruit thinning. Efficient harvesting of mature fruit often depends on some chemical promotion of abscission. Ethaphon and other compounds that can release ethylene are often used, although they sometimes show undesirable side effects of subsequent growth. For the harvesting of oranges in Florida, various combinations of cycloheximide, 5-chloro-3-methyl-4-nitro-1Hpyrazole, and chlorothalonil have been very effective in promoting abscission. See PLANT PHYSIOLOGY. Fredrick T. Addicott Bibliography. F. T. Addicott, Abscission, 1982; L. G. Nickell, Plant Growth Regulators, 1982.
Absolute zero The temperature at which an ideal gas would exert no pressure. The Kelvin scale of temperatures is defined in terms of the triple point of water, T3 = 273.16◦ (where the solid, liquid, and vapor phases coexist) and absolute zero. Thus, the Kelvin degree is 1/273.16 of the thermodynamic temperature of the triple point. Temperature is measured most simply via the constant-volume ideal-gas thermometer, in which a small amount of gas is introduced (in order to limit the effect of interactions between molecules) and then sealed off. The gas pressure P referenced to its value at the triple point P(T3) is measured, and
temperature T may be inferred from Eq. (1). The T (K) = 273.16 lim
P→0
P(T ) P(T3 )
(1)
ideal-gas law applies if the molecules in a gas exert no forces on one another and if they are not attracted to the walls. With helium gas, the constant-volume gas thermometer can be used down to a temperature on the order of 1 K. Absolute zero is the temperature at which the pressure of a truly ideal gas would vanish. See GAS THERMOMETRY; TEMPERATURE MEASUREMENT. In order that a temperature be assigned to a particular specimen, there must exist mechanisms for the free transfer of energy within the system. Measurement of the distribution of speeds of particles can serve to define the temperature through the Maxwell distribution. The most probable speed gets smaller as the temperature decreases. According to classical physics, all motion would cease at absolute zero; however, the quantum-mechanical uncertainty principle requires that there be a small amount of residual motion (zero-point motion) even at absolute zero. See KINETIC THEORY OF MATTER; QUANTUM MECHANICS; UNCERTAINTY PRINCIPLE. The Kelvin scale can be obtained by measurements of the efficiency of an ideal heat engine operating between reservoirs at two fixed temperatures. According to the second law of thermodynamics, the maximum possible efficiency is given by Eq. (2), =1−
T1 T2
(2)
where T1 and T2 are the temperatures of the lowand the high-temperature reservoirs respectively. All comparisons between the gas-thermometer scales and the thermodynamic scales have shown the equivalence of the two approaches. See CARNOT CYCLE; THERMODYNAMIC PRINCIPLES. Temperature can also be defined from the Boltzmann distribution. If a collection of spin 1/2 magnetic ions is placed in a magnetic field, the ratio of the occupancy of the lower to the higher energy state is given by Eq. (3). Here k is Boltzmann’s conNL |E| = exp NH kT
(3)
stant, |E| is the magnitude of the difference in energy between the states, and T is the Kelvin temperature. Thus, at high temperatures the two states have nearly equal occupation probability, while the lower energy state is progressively favored at lower temperatures. At absolute zero, only the lower energy level is occupied. By measuring the populations of the two states, the temperature of the spins can be inferred. This relation allows for the possibility of negative temperatures when the population of the higher energy state exceeds that of the lower state. From the point of view of energy content, negative temperatures correspond to an energy of the spin system that exceeds that of an infinite positive temperature, and thus they are hotter than ordinary
19
20
Absorption temperatures. See BOLTZMANN CONSTANT; BOLTZMANN STATISTICS; NEGATIVE TEMPERATURE. Negative temperatures notwithstanding, the third law of thermodynamics states that the absolute zero of temperature cannot be attained by any finite number of steps. The lowest (and hottest) temperatures that have been achieved are on the order of a picokelvin (10−12K). These are spin temperatures of nuclei which are out of equilibrium with the lattice vibrations and electrons of a solid. The lowest temperatures to which the electrons have been cooled are on the order of 10 microkelvins in metallic systems. See LOW-TEMPERATURE PHYSICS; TEMPERATURE. Jeevak M. Parpia; David M. Lee
Absorption Either the taking up of matter in bulk by other matter, as in the dissolving of a gas by a liquid; or the taking up of energy from radiation by the medium through which the radiation is passing. In the first case, an absorption coefficient is defined as the amount of gas dissolved at standard conditions by a unit volume of the solvent. Absorption in this sense is a volume effect: The absorbed substance permeates the whole of the absorber. In absorption of the second type, attenuation is produced which in many cases follows Lambert’s law and adds to the effects of scattering if the latter is present. See ATTENUATION; SCATTERING OF ELECTROMAGNETIC RADIATION. Energy. Absorption of electromagnetic radiation can occur in several ways. For example, microwaves in a waveguide lose energy to the walls of the guide. For nonperfect conductors, the wave penetrates the guide surface, and energy in the wave is transferred to the atoms of the guide. Light is absorbed by atoms of the medium through which it passes, and in some cases this absorption is quite distinctive. Selected frequencies from a heterochromatic source are strongly absorbed, as in the absorption spectrum of the Sun. Electromagnetic radiation can be absorbed by the photoelectric effect, where the light quantum is absorbed and an electron of the absorbing atom is ejected, and also by Compton scattering. Electronpositron pairs may be created by the absorption of a photon of sufficiently high energy. Photons can be absorbed by photoproduction of nuclear and subnuclear particles, analogous to the photoelectric effect. See ABSORPTION OF ELECTROMAGNETIC RADIATION; COMPTON EFFECT; ELECTRON-POSITRON PAIR PRODUCTION; PHOTOEMISSION. Sound waves are absorbed at suitable frequencies by particles suspended in the air (wavelength of the order of the particle size), where the sound energy is transformed into vibrational energy of the absorbing particles. See SOUND ABSORPTION. Absorption of energy from a beam of particles can occur by the ionization process, where an electron in the medium through which the beam passes is removed by the beam particles. The finite range of protons and alpha particles in matter is a result of this process. In the case of low-energy electrons,
scattering is as important as ionization, so that range is a less well-defined concept. Particles themselves may be absorbed from a beam. For example, in a nuclear reaction an incident particle X is absorbed into nucleus Y, and the result may be that another particle Z, or a photon, or particle X with changed energy comes out. Low-energy positrons are quickly absorbed by annihilating with electrons in matter to yield two gamma rays. See NUCLEAR REACTION. McAllister H. Hull, Jr. Matter. The absorption of matter is a chemical engineering unit operation. In the chemical process industries and in related areas such as petroleum refining and fuels purification, absorption usually means gas absorption. This is a unit operation in which a gas (or vapor) mixture is contacted with a liquid solvent selected to preferentially absorb one, or in some cases more than one, component from the mixture. The purpose is either to recover a desired component from a gas mixture or to rid the mixture of an impurity. In the latter case, the operation is often referred to as scrubbing. The gas, once absorbed in the liquid, is in a completely dissolved state and therefore has become a full-fledged member of the liquid phase. In the chemical industry, gas absorption is the second-most common separation operation involving gas-liquid contacting. Only fractional distillation is employed more frequently. When the operation is employed in reverse, that is, when a gas is utilized to extract a component from a liquid mixture, it is referred to as gas desorption, stripping, or sparging. In gas absorption, either no further changes occur to the gaseous component once it is absorbed in the liquid solvent, or the absorbed component (solute) will become involved in a chemical reaction with the solvent in the liquid phase. In the former case, the operation is referred to as physical gas absorption, and in the latter case as gas absorption with chemical reaction. Examples of simple, or physical, gas absorption include the absorption of light oil (benzene, toluene, xylenes) from coke oven by-product gases by scrubbing with a petroleum oil, and recovery of valuable solvent vapors, such as acetone or those used in drycleaning processes, from gas streams by washing the gas with an appropriate solvent for the vapors. Examples of gas absorption with chemical reaction include scrubbing the flue gases from coalfired generating stations with aqueous sodium carbonate solution to remove sulfur dioxide (a potential cause of acid rain), recovery of ammonia from coke oven by-product gases by scrubbing with dilute sulfuric acid, and the purification of natural gas by absorption of hydrogen sulfide in aqueous ethanolamine. Gas absorption is normally carried out in the same types of columns as those used in fractional distillation, namely tray-type or packed towers, but with the latter used much more frequently than the former in distillation. The operation can be carried out by using cocurrent flow, that is, with both the gas and liquid phases flowing upward in the column, but it is normally countercurrent with the liquid phase flowing
Absorption (biology) downward and the gas phase rising. Gas absorption, unlike distillation, is not carried out at the boiling point. To the contrary, and because the solubility of gases in liquids increases with decreasing temperature, gas absorption is usually carried out cold, that is, at a temperature not far above the freezing point of the solvent. An evolution of heat occurs during gas absorption, as a result of the absorbed component giving up its latent heat of vaporization upon absorption, just as it would have if it had been condensed. The resulting temperature rise is minimized by a normally large solvent-to-gas ratio used in order to keep low the concentration of absorbed gas in the liquid phase, and therefore to keep high the concentration driving force causing the absorption to occur. Unlike fractional distillation, gas absorption does not require a reboiler, condenser, or reflux. However, it does require a subsequent solvent-recovery operation to separate solvent and solute, so that the solvent can be recycled. The energy requirements for gas absorption and solvent recovery combined are usually substantially less than for fractional distillation. The main basis for solvent selection is that the solvent have a high solubility for the gaseous component to be absorbed, and the lowest possible solubility for the other components of the gas stream—in other words, that the solvent have both a large absorptive capacity and a high degree of selectivity. See GAS ABSORPTION OPERATIONS; DISTILLATION; DISTILLATION COLUMN; UNIT OPERATIONS. William F. Furter
Absorption (biology) The net movement (transport) of water and solutes from outside an organism to its interior. The unidirectional flow of materials into an animal from the environment generally takes place across the alimentary tract, the lungs, or the skin, and in each location a specific cell layer called an epithelium regulates the passage of materials. See EPITHELIUM; RESPIRATION. Structural aspects. The mammalian intestine is a good example of an epithelium which absorbs ions and water, along with a wide range of nutrient materials. The cells lining the gastrointestinal epithelium possess morphological features similar to those of many other epithelia. Most intestinal epithelia have cells which are columnar in shape, have an extensive microvillar membrane facing the interior (or lumen) of the gut, and are functionally specialized to transport molecules from the gut contents to the blood (Fig. 1). The brush border or microvillar apical membrane has an expanded surface to maximize uptake of materials into the cytoplasm. Adjacent cells are attached to one another at the apical surface (facing the lumen) by “tight junctions,” which serve as partial barriers to the passage of materials between the lumen and the blood by intercellular paths. Several types of enzymes, including peptidases, carbohydrases, and lipases, are located in the brush border membrane to break down complex peptides, carbohydrates, and fats before they are transported into the
tight junction intermediate junction
microvillus terminal web
desmosome microtubules mitochondria unattached ribosomes
lysosomes granular reticulum smooth-surfaced reticulum
Golgi material intercellular space
nucleus
basement membrane lamina propria
Fig. 1. Diagram of an absorptive epithelial cell from the mammalian small intestine. (After L. R. Johnson, ed., Physiology of the Gastrointestinal Tract, Raven Press, 1981)
epithelial cytoplasm as molecular monomers (single units). See DIGESTIVE SYSTEM. The sides and bottoms of the cells with respect to the lumen are called the lateral and basal surfaces, or sometimes collectively the basolateral surfaces, and these may vary considerably among different epithelia and among species. While some epithelial cells have long, straight basolateral surfaces, others may be highly interdigitated with neighboring cells to increase their total surface area. In some epithelia these basolateral interdigitations contain abundant mitochondria, which provide energy for transport of materials across membranes. See MITOCHONDRIA. During epithelial absorption, molecules must first pass across the apical membrane (from the mucosal side) into the cytoplasm. This entry process is then followed by ejection of the same molecule across the basolateral border to the blood side (called the serosal side) for eventual distribution to the body via the circulatory system. See CELL MEMBRANES. Kinds of transport. Absorption across epithelia may occur by several different passive and active processes. Simple diffusion is the net movement of molecules from the apical to basolateral surfaces of an epithelium down chemical and electrical gradients without the requirement of cellular energy sources. Facilitated diffusion across the epithelium is similar to simple diffusion in that energy is not required, but in this process, molecular interaction with protein binding sites (carriers) in one or both membranes must occur to facilitate the transfer. Active molecular transport involves the use of membrane protein carriers as well as cellular energy supplies to move a transported molecule up an electrochemical gradient across the epithelium. Endocytosis and phagocytosis are also examples of active
21
22
Absorption (biology) transport because metabolic energy is required, but in these processes whole regions of the cell membrane are used to engulf fluid or particles, rather than to bring about molecular transfer using singlemembrane proteins. See ENDOCYTOSIS; PHAGOCYTOSIS. Mechanisms of ion absorption. Although a wide variety of ions are absorbed by different types of epithelial cells, the mechanisms of Na+ and Cl− transport in mammalian small intestine are perhaps best known in detail. Transepithelial transport of these two ions occurs in this tissue by three independent processes: active Na+ absorption, not coupled directly to the flow of other solutes but accompanied indirectly by the diffusional absorption of Cl−; coupled NaCl absorption; and cotransport of Na+ with a wide variety of nutrient molecules. During uncoupled active Na+ transport, the cation crosses the apical membrane of the small-intestinal epithelial cell down an electrochemical gradient. Transfer across the brush border may be by simple or facilitated diffusion. Efflux of Na+ up a concentration gradient out of the cell across the basolateral surface occurs by active transport, using a carrier protein called the Na+/K+–adenosine triphosphatase (Na+ pump). Absorption of Na+ by this two-step process results in the generation of a transepithelial electrical potential difference between the lumen and the blood of about 3–5 millivolts (blood electrically positive to apical surface). This electrical potential difference passively attracts Cl− from the luminal side, most of this diffusion occurring paracellularly between the epithelial cells. Two mechanisms have been proposed for coupled NaCl absorption in the mammalian small intestine. In the first model (Fig. 2a), Na+ and Cl− bind to common apical membrane carriers and enter the epithelium together, with Na+ diffusing down an electrochemical gradient and Cl− moving up an electro-
Mucosal solution
Serosal solution
Cell CI
Na
Na cAMP CI
CI
?
Na
(a)
HCO3
H2O + CO2
CI
CO2
H2CO3 Na (b)
H
Fig. 2. Two models for electroneutral NaCl absorption by mammalian small-intestinal epithelial cells from the mucosal solution (food side) to the serosal solution (blood side). (a) Ions enter together by binding to common apical membrane carriers. (b) Ions enter by separate brush border carriers. (After L. R. Johnson, ed., Physiology of the Gastrointestinal Tract, Raven Press, 1981)
water
water flow solute pumps
Fig. 3. Mechanism of solute-linked water absorption by epithelial cells. Density of dots refers to relative osmotic pressure in different compartments. (After R. Eckert and D. Randall, Animal Physiology, Freeman, 1983)
chemical gradient. Efflux of Na+ to the blood occurs by way of the energy-dependent Na+ pump, while Cl− moves passively across the basolateral border by facilitated diffusion. In the second model (Fig. 2b), Na+ and Cl− enter the cell by separate brush border carriers in exchange for intracellular H+ and HCO3−, respectively. Their efflux to the blood occurs as described in Fig. 2a. A final mechanism leading to the net flow of Na+ across an epithelium is by coupling the entry of this ion into the cell with that of a wide variety of organic solutes. This pathway of Na+ transport will be discussed below in conjunction with the mechanisms of nutrient absorption. Mechanisms of water absorption. Net water transport across an epithelium is coupled to net ion transport in the same direction. As shown in Fig. 3, Na+pump sites are believed to be located along the lateral borders of epithelial cells. Energy-dependent Na+ efflux from the cells to the intercellular spaces creates a local increase in osmotic pressure within these small compartments. An osmotic pressure gradient becomes established here, with greatest solute concentrations located nearest the tight junctions. Water flows into the cell across the brush border membrane and out the lateral membranes in response to the increased osmotic pressure in the paracellular spaces. Once water is in the intercellular compartment, a buildup of hydrostatic pressure forces the transported fluid to the capillary network. See OSMOREGULATORY MECHANISMS. Mechanisms of nutrient absorption. It is believed that most nutrient molecules enter intestinal epithelial cells in conjunction with Na+ on a common brush border membrane carrier protein (Fig. 4). Several types of proteins may exist on this membrane for each major category of organic solute (for example, a protein for sugars is functionally distinct from one for amino acids), but each carrier type requires Na+ as a cotransported substrate. In all cases, Na+ enters the cell down an electrochemical gradient, while the nutrient molecule is accumulated against a concentration gradient within the cytoplasm. The diffusional permeability of the brush border membrane toward organic solutes is very low, so that even though a considerable concentration gradient for the nutrient builds up across the apical cell pole, little passive loss
Absorption of electromagnetic radiation 0 mV Na
− 40 mV
solute K
+ 5 mV Na K
solute
Na CI
solute
Fig. 4. Model for sodium-coupled absorption of organic solutes (D-hexoses, L-amino acids, dipeptides, tripeptides, water-soluble vitamins, bile salts) by mammalian small-intestinal epithelial cells. Cl− ions diffuse between epithelial cells, down the electrochemical gradient. (After L. R. Johnson, ed., Physiology of the Gastrointestinal Tract, Raven Press, 1981)
to the lumen takes place. Intracellular Na+ concentration is kept very low due to energy-dependent basolateral Na+-pump activities transferring this cation up a concentration gradient out of the cell. Accumulated intracellular nutrients are transferred across the basolateral border by facilitated diffusion on a carrier protein which does not have a Na+ requirement. No cellular metabolic energy is used directly for the absorption of nutrients. Instead, energy resources are used to maintain a steep Na+ gradient across the brush border membrane to facilitate the accumulation of nutrient across this cell pole as a result of its cotransport with the cation. See ION TRANSPORT. Gregory A. Ahearn Bibliography. G. Giebisch, D. C. Tosteson, and H. H. Ussing (eds.), Membrane Transport in Biology, vols. 3 and 4, 1978; L. R. Johnson (ed.), Physiology of the Gastrointestinal Tract, vols. 1 and 2, 1981; H. Murer, U. Hopfer, and R. Kinne, Sodiumproton antiport in brush border membrane vesicles isolated from rat small intestine, Biochem. J., 154:597–604, 1976; H. N. Nellans, R. A. Frizzell, and S. G. Schultz, Coupled sodium-chloride influx across the brush border of rabbit ileum, Amer. J. Physiol., 225:467–475, 1973; S. G. Schultz and P. F. Curran, Coupled transport of sodium and organic solutes, Physiol. Rev., 50:637–718, 1970.
Absorption of electromagnetic radiation The process whereby the intensity of a beam of electromagnetic radiation is attenuated in passing through a material medium by conversion of the energy of the radiation to an equivalent amount of energy which appears within the medium; the radiant energy is converted into heat or some other form of molecular energy. A perfectly transparent medium permits the passage of a beam of radiation without any change in intensity other than that caused by the spread or convergence of the beam, and the total radiant energy emergent from such a medium equals that which entered it, whereas the emergent energy from an absorbing medium is less than that which
enters, and, in the case of highly opaque media, is reduced practically to zero. No known medium is opaque to all wavelengths of the electromagnetic spectrum, which extends from radio waves, whose wavelengths are measured in kilometers, through the infrared, visible, and ultraviolet spectral regions, to x-rays and gamma rays, of wavelengths down to 10−13 m. Similarly, no material medium is transparent to the whole electromagnetic spectrum. A medium which absorbs a relatively wide range of wavelengths is said to exhibit general absorption, while a medium which absorbs only restricted wavelength regions of no great range exhibits selective absorption for those particular spectral regions. For example, the substance pitch shows general absorption for the visible region of the spectrum, but is relatively transparent to infrared radiation of long wavelength. Ordinary window glass is transparent to visible light, but shows general absorption for ultraviolet radiation of wavelengths below about 310 nanometers, while colored glasses show selective absorption for specific regions of the visible spectrum. The color of objects which are not self-luminous and which are seen by light reflected or transmitted by the object is usually the result of selective absorption of portions of the visible spectrum. Many colorless substances, such as benzene and similar hydrocarbons, selectively absorb within the ultraviolet region of the spectrum, as well as in the infrared. See COLOR; ELECTROMAGNETIC RADIATION. Laws of absorption. The capacity of a medium to absorb radiation depends on a number of factors, mainly the electronic and nuclear constitution of the atoms and molecules of the medium, the wavelength of the radiation, the thickness of the absorbing layer, and the variables which determine the state of the medium, of which the most important are the temperature and the concentration of the absorbing agent. In special cases, absorption may be influenced by electric or magnetic fields. The state of polarization of the radiation influences the absorption of media containing certain oriented structures, such as crystals of other than cubic symmetry. See STARK EFFECT; ZEEMAN EFFECT. Lambert’s law. Lambert’s law, also called Bouguer’s law or the Lambert-Bouguer law, expresses the effect of the thickness of the absorbing medium on the absorption. If a homogeneous medium is thought of as being constituted of layers of uniform thickness set normally to the beam, each layer absorbs the same fraction of radiation incident on it. If I is the intensity to which a monochromatic parallel beam is attenuated after traversing a thickness d of the medium, and I0 is the intensity of the beam at the surface of incidence (corrected for loss by reflection from this surface), the variation of intensity throughout the medium is expressed by Eq. (1), in which α is a I = I0 e−αd
(1)
constant for the medium called the absorption coefficient. This exponential relation can be expressed in
23
24
Absorption of electromagnetic radiation an equivalent logarithmic form as in Eq. (2), where log10 (I0 /I) = (α/2.303)d = kd
(2)
k = α/2.303 is called the extinction coefficient for radiation of the wavelength considered. The quantity log10 (I0/I) is often called the optical density, or the absorbance of the medium. Equation (2) shows that as monochromatic radiation penetrates the medium, the logarithm of the intensity decreases in direct proportion to the thickness of the layer traversed. If experimental values for the intensity of the light emerging from layers of the medium of different thicknesses are available (corrected for reflection losses at all reflecting surfaces), the value of the extinction coefficient can be readily computed from the slope of the straight line representing the logarithms of the emergent intensities as functions of the thickness of the layer. Equations (1) and (2) show that the absorption and extinction coefficients have the dimensions of reciprocal length. The extinction coefficient is equal to the reciprocal of the thickness of the absorbing layer required to reduce the intensity to one-tenth of its incident value. Similarly, the absorption coefficient is the reciprocal of the thickness required to reduce the intensity to 1/e of the incident value, where e is the base of the natural logarithms, 2.718. Beer’s law. This law refers to the effect of the concentration of the absorbing medium, that is, the mass of absorbing material per unit of volume, on the absorption. This relation is of prime importance in describing the absorption of solutions of an absorbing solute, since the solute’s concentration may be varied over wide limits, or the absorption of gases, the concentration of which depends on the pressure. According to Beer’s law, each individual molecule of the absorbing material absorbs the same fraction of the radiation incident upon it, no matter whether the molecules are closely packed in a concentrated solution or highly dispersed in a dilute solution. The relation between the intensity of a parallel monochromatic beam which emerges from a plane parallel layer of absorbing solution of constant thickness and the concentration of the solution is an exponential one, of the same form as the relation between intensity and thickness expressed by Lambert’s law. The effects of thickness d and concentration c on absorption of monochromatic radiation can therefore be combined in a single mathematical expression, given in Eq. (3), in which k is a constant for a
I = I0 e−k cd
(3)
given absorbing substance (at constant wavelength and temperature), independent of the actual concentration of solute in the solution. In logarithms, the relation becomes Eq. (4). The values of the constants k log10 (I0 /I) = (k /2.303)cd = cd
(4)
and in Eqs. (3) and (4) depend on the units of concentration. If the concentration of the solute is expressed in moles per liter, the constant is called the
molar extinction coefficient. Some authors employ the symbol aM, which is called the molar absorbance index, instead of . If Beer’s law is adhered to, the molar extinction coefficient does not depend on the concentration of the absorbing solute, but usually changes with the wavelength of the radiation, with the temperature of the solution, and with the solvent. The dimensions of the molar extinction coefficient are reciprocal concentration multiplied by reciprocal length, the usual units being liters/(mole)(cm). If Beer’s law is true for a particular solution, the plot of log (I0/I) against the concentrations for solutions of different concentrations, measured in cells of constant thickness, will yield a straight line, the slope of which is equal to the molar extinction coefficient. While no true exceptions to Lambert’s law are known, exceptions to Beer’s law are not uncommon. Such exceptions arise whenever the molecular state of the absorbing solute depends on the concentration. For example, in solutions of weak electrolytes, whose ions and undissociated molecules absorb radiation differently, the changing ratio between ions and undissociated molecules brought about by changes in the total concentration prevents solutions of the electrolyte from obeying Beer’s law. Aqueous solutions of dyes frequently deviate from the law because of dimerization and more complicated aggregate formation as the concentration of dye is increased. Absorption measurement. The measurement of the absorption of homogeneous media is usually accomplished by absolute or comparative measurements of the intensities of the incident and transmitted beams, with corrections for any loss of radiant energy caused by processes other than absorption. The most important of these losses is by reflection at the various surfaces of the absorbing layer and of vessels which may contain the medium, if the medium is liquid or gaseous. Such losses are usually automatically compensated for by the method of measurement employed. Losses by reflection not compensated for in this manner may be computed from Fresnel’s laws of reflection. See REFLECTION OF ELECTROMAGNETIC RADIATION. Scattering. Absorption of electromagnetic radiation should be distinguished from the phenomenon of scattering, which occurs during the passage of radiation through inhomogeneous media. Radiant energy which traverses media constituted of small regions of refractive index different from that of the rest of the medium is diverted laterally from the direction of the incident beam. The diverted radiation gives rise to the hazy or opalescent appearance characteristic of such media, exemplified by smoke, mist, and opal. If the centers of inhomogeneity are sufficiently dilute, the intensity of a parallel beam is diminished in its passage through the medium because of the sidewise scattering, according to a law of the same form as the Lambert-Bouguer law for absorption, given in Eq. (5), where I is the intensity of I = I0 e−τ d
(5)
Absorption of electromagnetic radiation the primary beam of initial intensity I0, after it has traversed a distance d through the scattering medium. The coefficient τ , called the turbidity of the medium, plays the same part in weakening the primary beam by scattering as does the absorption coefficient in true absorption. However, in true scattering, no loss of total radiant energy takes place, energy lost in the direction of the primary beam appearing in the radiation scattered in other directions. In some inhomogeneous media, both absorption and scattering occur together. See SCATTERING OF ELECTROMAGNETIC RADIATION. Physical nature. Absorption of radiation by matter always involves the loss of energy by the radiation and a corresponding gain in energy by the atoms or molecules of the medium. The energy of an assembly of gaseous atoms consists partly of kinetic energy of the translational motion which determines the temperature of the gas (thermal energy), and partly of internal energy, associated with the binding of the extranuclear electrons to the nucleus, and with the binding of the particles within the nucleus itself. Molecules, composed of more than one atom, have, in addition, energy associated with periodic rotations of the molecule as a whole and with oscillations of the atoms within the molecule with respect to one another. The energy absorbed from radiation appears as increased internal energy, or in increased vibrational and rotational energy of the atoms and molecules of the absorbing medium. As a general rule, translational energy is not directly increased by absorption of radiation, although it may be indirectly increased by degradation of electronic energy or by conversion of rotational or vibrational energy to that of translation by intermolecular collisions. Quantum theory. In order to construct an adequate theoretical description of the energy relations between matter and radiation, it has been necessary to amplify the wave theory of radiation by the quantum theory, according to which the energy in radiation occurs in natural units called quanta. The value of the energy in these units, expressed in ergs or calories, for example, is the same for all radiation of the same wavelength, but differs for radiation of different wavelengths. The energy E in a quantum of radiation of frequency ν (where the frequency is equal to the velocity of the radiation in a given medium divided by its wavelength in the same medium) is directly proportional to the frequency, or inversely proportional to the wavelength, according to the relation given in Eq. (6), where h is a universal conE = hν
(6)
stant known as Planck’s constant. The value of h is 6.63 × 10−34 joule-second, and if ν is expressed in s−1, E is given in joules per quantum. See QUANTUM MECHANICS. The most energetic type of change that can occur in an atom involves the nucleus, and increase of nuclear energy by absorption therefore requires quanta of very high energy, that is, of high frequency or low
wavelength. Such rays are the γ -rays, whose wavelength varies downward from 0.01 nm. Next in energy are the electrons nearest to the nucleus and therefore the most tightly bound. These electrons can be excited to states of higher energy by absorption of x-rays, whose range in wavelength is from about 0.01 to 1 nm. Less energy is required to excite the more loosely bound valence electrons. Such excitation can be accomplished by the absorption of quanta of visible radiation (wavelength 700 nm for red light to 400 nm for blue) or of ultraviolet radiation, of wavelength down to about 100 nm. Absorption of ultraviolet radiation of shorter wavelengths, down to those on the border of the x-ray region, excites electrons bound to the nucleus with intermediate strength. The absorption of relatively low-energy quanta of wavelength from about 1 to 10 micrometers suffices to excite vibrating atoms in molecules to higher vibrational states, while changes in rotational energy, which are of still smaller magnitude, may be excited by absorption of radiation of still longer wavelength, from the short-wavelength radio region of about 1 cm to long-wavelength infrared radiation, some hundredths of a centimeter in wavelength. Gases. The absorption of gases composed of atoms is usually very selective. For example, monatomic sodium vapor absorbs very strongly over two narrow wavelength regions in the yellow part of the visible spectrum (the so-called D lines), and no further absorption by monatomic sodium vapor occurs until similar narrow lines appear in the near-ultraviolet. The valence electron of the sodium atom can exist only in one of a series of energy states separated by relatively large energy intervals between the permitted values, and the sharp-line absorption spectrum results from transitions of the valence electron from the lowest energy which it may possess in the atom to various excited levels. Line absorption spectra are characteristic of monatomic gases in general. See ATOMIC STRUCTURE AND SPECTRA. The visible and ultraviolet absorption of vapors composed of diatomic or polyatomic molecules is much more complicated than that of atoms. As for atoms, the absorbed energy is utilized mainly in raising one of the more loosely bound electrons to a state of higher energy, but the electronic excitation of a molecule is almost always accompanied by simultaneous excitation of many modes of vibration of the atoms within the molecule and of rotation of the molecule as a whole. As a result, the absorption, which for an atom is concentrated in a very sharp absorption line, becomes spread over a considerable spectral region, often in the form of bands. Each band corresponds to excitation of a specific mode of vibration accompanying the electronic change, and each band may be composed of a number of very fine lines close together in wavelength, each of which corresponds to a specific rotational change of the molecule accompanying the electronic and vibrational changes. Band spectra are as characteristic of the absorption of molecules in the gaseous state, and frequently in the liquid state, as line spectra
25
26
Absorption of electromagnetic radiation are of gaseous atoms. See BAND SPECTRUM; LINE SPECTRUM; MOLECULAR STRUCTURE AND SPECTRA. Liquids. Liquids usually absorb radiation in the same general spectral region as the corresponding vapors. For example, liquid water, like water vapor, absorbs infrared radiation strongly (vibrational transitions), is largely transparent to visible and near-ultraviolet radiation, and begins to absorb strongly in the farultraviolet. A universal difference between liquids and gases is the disturbance in the energy states of the molecules in a liquid caused by the great number of intermolecular collisions; this has the effect of broadening the very fine lines observed in the absorption spectra of vapors, so that sharp-line structure disappears in the absorption bands of liquids. Solids. Substances which can exist in solid, liquid, and vapor states without undergoing a temperature rise to very high values usually absorb in the same general spectral regions for all three states of aggregation, with differences in detail because of the intermolecular forces present in the liquid and solid. Crystalline solids, such as rock salt or silver chloride, absorb infrared radiation of long wavelength, which excites vibrations of the electrically charged ions of which these salts are composed; such solids are transparent to infrared radiations of shorter wavelengths. In colorless solids, the valence electrons are too tightly bound to the nuclei to be excited by visible radiation, but all solids absorb in the near- or farultraviolet region. See CRYSTAL OPTICS; INTERMOLECULAR FORCES. The use of solids as components of optical instruments is restricted by the spectral regions to which they are transparent. Crown glass, while showing excellent transparency for visible light and for ultraviolet radiation immediately adjoining the visible region, becomes opaque to radiation of wavelength about 300 nm and shorter, and is also opaque to infrared radiation longer than about 2000 nm in wavelength. Quartz is transparent down to wavelengths about 180 nm in the ultraviolet, and to about 4 µm in the infrared. The most generally useful material for prisms and windows for the near-infrared region is rock salt, which is highly transparent out to about 15 µm. For a detailed discussion of the properties of optical glass see OPTICAL MATERIALS. Fluorescence. The energy acquired by matter by absorption of visible or ultraviolet radiation, although primarily used to excite electrons to higher energy states, usually ultimately appears as increased kinetic energy of the molecules, that is, as heat. It may, however, under special circumstances, be reemitted as electromagnetic radiation. Fluorescence is the reemission, as radiant energy, of absorbed radiant energy, normally at wavelengths the same as or longer than those absorbed. The reemission, as ordinarily observed, ceases immediately when the exciting radiation is shut off. Refined measurements show that the fluorescent reemission persists, in different cases, for periods of the order of 10−9 to 10−5 s. The simplest case of fluorescence is the resonance fluorescence of monatomic gases at low pressure, such as sodium or mercury vapors, in
which the reemitted radiation is of the same wavelength as that absorbed. In this case, fluorescence is the converse of absorption: Absorption involves the excitation of an electron from its lowest energy state to a higher energy state by radiation, while fluorescence is produced by the return of the excited electron to the lower state, with the emission of the energy difference between the two states as radiation. The fluorescent radiation of molecular gases and of nearly all liquids, solids, and solutions contains a large component of wavelengths longer than those of the absorbed radiation, a relationship known as Stokes’ law of fluorescence. In these cases, not all of the absorbed energy is reradiated, a portion remaining as heat in the absorbing material. The fluorescence of iodine vapor is easily seen on projecting an intense beam of visible light through an evacuated bulb containing a few crystals of iodine, but the most familiar examples are provided by certain organic compounds in solution—for instance, quinine sulfate, which absorbs ultraviolet radiation and reemits blue, or fluorescein, which absorbs blue-green light and fluoresces with an intense, bright-green color. See FLUORESCENCE. Phosphorescence. The radiant reemission of absorbed radiant energy at wavelengths longer than those absorbed, for a readily observable interval after withdrawal of the exciting radiation, is called phosphorescence. The interval of persistence, determined by means of a phosphoroscope, usually varies from about 0.001 s to several seconds, but some phosphors may be induced to phosphorescence by heat days or months after the exciting absorption. An important and useful class of phosphors is the impurity phosphors, solids such as the sulfides of zinc or calcium which are activated to the phosphorescent state by incorporating minute amounts of foreign material (called activators), such as salts of manganese or silver. So-called fluorescent lamps contain a coating of impurity phosphor on their inner wall which, after absorbing ultraviolet radiation produced by passage of an electrical discharge through mercury vapor in the lamp, reemits visible light. The receiving screen of a television tube contains a similar coating, excited not by radiant energy but by the impact of a stream of electrons on the surface. See FLUORESCENT LAMP; PHOSPHORESCENCE. Luminescence. Phosphorescence and fluorescence are special cases of luminescence, which is defined as light emission that cannot be attributed merely to the temperature of the emitting body. Luminescence may be excited by heat (thermoluminescence), by electricity (electroluminescence), by chemical reaction (chemiluminescence), or by friction (triboluminescence), as well as by radiation. See LUMINESCENCE. Absorption and emission coefficients. The absorption and emission processes of atoms were examined from the quantum point of view by Albert Einstein in 1916, with some important results that have been realized practically in the invention of the maser and the laser. Consider an assembly of atoms undergoing absorption transitions of frequency ν s−1 from the
Absorption of electromagnetic radiation ground state 1 to an excited state 2 and emission transitions in the reverse direction, the atoms and radiation being at equilibrium at temperature T. The equilibrium between the excited and unexcited atoms is determined by the Boltzmann relation N2/N1 = exp(−hν/kT), where N1 and N2 are the equilibrium numbers of atoms in states 1 and 2, respectively, and the radiational equilibrium is determined by equality in the rate of absorption and emission of quanta. The number of quanta absorbed per second is B12N1ρ(ν), where ρ(ν) is the density of radiation of frequency ν (proportional to the intensity), and B12 is a proportionality constant called the Einstein coefficient for absorption. Atoms in state 2 will emit radiation spontaneously (fluorescence), after a certain mean life, at a rate of A21N2 per second, where A21 is the Einstein coefficient for spontaneous emission from state 2 to state 1. To achieve consistency between the density of radiation of frequency ν at equilibrium calculated from these considerations and the value calculated from Planck’s radiation law, which is experimentally true, it is necessary to introduce, in addition to the spontaneous emission, an emission of intensity proportional to the radiation density of frequency ν in which the atoms are immersed. The radiational equilibrium is then determined by Eq. (7), where B21 is
Dispersion. A transparent material does not abstract energy from radiation which it transmits, but it always decreases the velocity of propagation of such radiation. In a vacuum, the velocity of radiation is the same for all wavelengths, but in a material medium, the velocity of propagation varies considerably with wavelength. The refractive index n of a medium is the ratio of the velocity of light in vacuum to that in the medium, and the effect of the medium on the velocity of radiation which it transmits is expressed by the variation of refractive index with the wavelength λ of the radiation, dn/dλ. This variation is called the dispersion of the medium. For radiation of wavelengths far removed from those of absorption bands of the medium, the refractive index increases regularly with decreasing wavelength or increasing frequency; the dispersion is then said to be normal. In regions of normal dispersion, the variation of refractive index with wavelength can be expressed with considerable accuracy by Eq. (10), known as
B12 N1 ρ(ν) = A21 N2 + B21 N2 ρ(ν)
dn −2B = 3 dλ λ
(7)
the Einstein coefficient of stimulated emission. The Einstein radiation coefficients are found to be related by Eqs. (8a) and (8b). B12 = B21
(8a)
A21 = (8πhν 3 /c3 ) · B21
(8b)
In the past when one considered radiation intensities available from terrestrial sources, stimulated emission was very feeble compared with the spontaneous process. Stimulated emission is, however, the fundamental emission process in the laser, a device in which a high concentration of excited molecules is produced by intense illumination from a “pumping” source, in an optical system in which excitation and emission are augmented by back-and-forth reflection until stimulated emission swamps the spontaneous process. See LASER; OPTICAL PUMPING. There are also important relations between the absorption characteristics of atoms and their mean lifetime τ in the excited state. Since A21 is the number of times per second that a given atom will emit a quantum spontaneously, the mean lifetime before emission in the excited state is τ = 1/A21. It can also be shown that A21 and τ are related, as shown in Eq. (9), (8π 2 ν 2 e2 ) ·f mc3 = 7.42 × 10−22 f ν 2 (ν in s−1 )
A21 = 1/τ =
(9)
to the f number or oscillator strength for the transition that occurs in the dispersion equations shown as Eqs. (13) to (17). The value of f can be calculated from the absorption integrated over the band according to Eq. (18).
n=A+
B C + 4 λ2 λ
(10)
Cauchy’s equation, in which A, B, and C are constants with positive values. As an approximation, C may be neglected in comparison with A and B, and the dispersion, dn/dλ, is then given by Eq. (11). Thus, in (11)
regions of normal dispersion, the dispersion is approximately inversely proportional to the cube of the wavelength. Dispersion by a prism. The refraction, or bending, of a ray of light which enters a material medium obliquely from vacuum or air (the refractive index of which for visible light is nearly unity) is the result of the diminished rate of advance of the wavefronts in the medium. Since, if the dispersion is normal, the refractive index of the medium is greater for violet than for red light, the wavefront of the violet light is retarded more than that of the red light. Hence, white light entering obliquely into the medium is converted within the medium to a continuously colored band, of which the red is least deviated from the direction of the incident beam, the violet most, with orange, yellow, green, and blue occupying intermediate positions. On emergence of the beam into air again, the colors remain separated. The action of the prism in resolving white light into its constituent colors is called color dispersion. See OPTICAL PRISM; REFRACTION OF WAVES. The angular dispersion of a prism is the ratio, dθ/dλ, of the difference in angular deviation dθ of two rays of slightly different wavelength which pass through the prism to the difference in wavelength dλ when the prism is set for minimum deviation. The angular dispersion of the prism given in Eq. (12) is the product of two factors, the variation, dθ dθ dn = · dλ dn dλ
(12)
dθ/dn, the deviation θ with refractive index n, and
27
28
Absorption of electromagnetic radiation the variation of refractive index with wavelength, the dispersion of the material of which the prism is made. The latter depends solely on this material, while dθ /dn depends on the angle of incidence and the refracting angle of the prism. The greater the dispersion of the material of the prism, the greater is the angular separation between rays of two given wavelengths as they leave the prism. For example, the dispersion of quartz for visible light is lower than that of glass; hence the length of the spectrum from red to violet formed by a quartz prism is less than that formed by a glass prism of equal size and shape. Also, since the dispersion of colorless materials such as glass or quartz is greater for blue and violet light than for red, the red end of the spectrum formed by prisms is much more contracted than the blue. The colors of the rainbow result from dispersion of sunlight which enters raindrops and is refracted and dispersed in passing through them to the rear surface, at which the dispersed rays are reflected and reenter the air on the side of the drop on which the light was incident. See METEOROLOGICAL OPTICS; RAINBOW. Anomalous dispersion. The regular increase of refractive index with decreasing wavelength expressed by Cauchy’s equation breaks down as the wavelengths approach those of strong absorption bands. As the absorption band is approached from the longwavelength side, the refractive index becomes very large, then decreases within the band to assume abnormally small values on the short-wavelength side, values below those for radiation on the longwavelength side. A hollow prism containing an alcoholic solution of the dye fuchsin, which absorbs green light strongly, forms a spectrum in which the violet rays are less deviated than the red, on account of the abnormally low refractive index of the medium for violet light. The dispersion of media for radiation of wavelengths near those of strong absorption bands is said to be anomalous, in the sense that the refractive index decreases with decreasing wavelength instead of showing the normal increase. The theory of dispersion shows, however, that both the normal and anomalous variation of refractive index with wavelength can be satisfactorily described as aspects of a unified phenomenon, so that there is nothing fundamentally anomalous about dispersion in the vicinity of an absorption band. See DISPERSION (RADIATION). Normal and anomalous dispersion of quartz are illustrated in Fig. 1. Throughout the near-infrared, visible, and near-ultraviolet spectral regions (between P and R on the curve), the dispersion is normal and adheres closely to Cauchy’s equation, but it becomes anomalous to the right of R. From S to T, Cauchy’s equation is again valid. Relation to absorption. Figure 1 shows there is an intimate connection between dispersion and absorption; the refractive index rises to high values as the absorption band is approached from the long– wavelength side and falls to low values on the shortwavelength side of the band. In fact, the theory of dispersion shows that the complete dispersion curve
n
P Q R
Cauchy equation
S T
A
A
visible region
absorption band
λ
Fig. 1. Curve showing anomalous dispersion of quartz. A is limiting value of n as λ approaches infinity. (After F. A. Jenkins and H. E. White, Fundamentals of Optics, 4th ed., McGraw-Hill, 1976)
as a function of wavelength is governed by the absorption bands of the medium. In classical electromagnetic theory, electric charges are regarded as oscillating, each with its appropriate natural frequency ν 0, about positions of equilibrium within atoms or molecules. Placed in a radiation field of frequency ν per second, the oscillator in the atom is set into forced vibration, with the same frequency as that of the radiation. When ν is much lower or higher than ν 0, the amplitude of the forced vibration is small, but the amplitude becomes large when the frequency of the radiation equals the natural frequency of the oscillator. In much the same way, a tuning fork is set into vibration by sound waves corresponding to the same note emitted by another fork vibrating at the same frequency. To account for the absorption of energy by the medium from the radiation, it is necessary to postulate that in the motion of the atomic oscillator some frictional force, proportional to the velocity of the oscillator, must be overcome. For small amplitudes of forced oscillation, when the frequency of the radiation is very different from the natural period of the oscillator, the frictional force and the absorption of energy are negligible. Near resonance between the radiation and the oscillator, the amplitude becomes large, with a correspondingly large absorption of energy to overcome the frictional resistance. Radiation of frequencies near the natural frequency therefore corresponds to an absorption band. See SYMPATHETIC VIBRATION. To show that the velocity of the radiation within the medium is changed, it is necessary to consider the phase of the forced vibration, which the theory shows to depend on the frequency of the radiation. The oscillator itself becomes a source of secondary radiation waves within the medium which combine to form sets of waves moving parallel to the original waves. Interference between the secondary and primary waves takes place, and because the phase of the secondary waves, which is the same as that of the atomic oscillators, is not the same as that of the primary waves, the wave motion resulting from the interference between the two sets of waves is different in phase from that of the primary waves incident
Absorption of electromagnetic radiation on the medium. But the velocity of propagation of the waves is the rate of advance of equal phase; hence the phase change effected by the medium, which is different for each frequency of radiation, is equivalent to a change in the velocity of the radiation within the medium. When the frequency of the radiation slightly exceeds the natural frequency of the oscillator, the radiation and the oscillator becomes 180◦ out of the phase, which corresponds to an increase in the velocity of the radiation and accounts for the fall in refractive index on the short-wavelength side of the absorption band. The theory leads to Eqs. (13) through (17) for the refractive index of a material medium as a function of the frequency of the radiation. In the equations the frequency is expressed as angular frequency, ω = 2π ν s−1 = 2π c/λ where c is the velocity of light. When the angular frequency ω of the radiation is not very near the characteristic frequency of the electronic oscillator, the refractive index of the homogeneous medium containing N molecules per cubic centimeter is given by Eq. (13a), where e and m are f 4π Ne2 · 2 m ω0 − ω 2 fi /mi n2 = 1 + 4πNe2 ωi 2 − ω 2 i n2 = 1 +
(13a) (13b)
the charge and mass of the electron, and f is the number of the oscillators per molecule of characteristics frequency ω0. The f value is sometimes called the oscillator strength. If the molecule contains oscillators of different frequencies and mass (for example, electronic oscillators of frequency corresponding to ultraviolet radiation and ionic oscillators corresponding to infrared radiation), the frequency term becomes a summation, as in Eq. (13b), where ωi is the characteristic frequency of the ith type of oscillator, and fi and mi are the corresponding f value and mass. In terms of wavelengths, this relation can be written as Eq. (14), where Ai is a constant for the medium, n2 = 1 +
i
Ai λ2 − λi 2
λ2
(14)
λ is the wavelength of the radiation, and λi = c/ν i is the wavelength corresponding to the characteristic frequency ν i per second (Sellmeier’s equation). If the medium is a gas, for which the refractive index is only slightly greater than unity, the dispersion formula can be written as Eq. (15). n = 1 + 2π Ne2
fi /mi ωi 2 − ω 2 i
The absorption of radiant energy of frequency very close to the characteristic frequency of the medium is formally regarded as the overcoming of a frictional force when the molecular oscillators are set into vibration, related by a proportionality constant g to the velocity of the oscillating particle; g is a damping coefficient for the oscillation. If the refractive index is determined by a single electronic oscillator, the dispersion equation for a gas at radiational frequencies within the absorption band becomes Eq. (16). n=1+
2πNe2 f (ω0 2 − ω2 ) 2 m (ω0 − ω2 )2 + ω2 g2
(16)
At the same time an absorption constant κ enters the equations, related to the absorption coefficient α of Eq. (1) by the expression κ = αc/2ωµ. Equation (17) shows the relationship. For a monatomic κ=
2πNe2 f ωg m (ω0 2 − ω2 )2 + ω2 g2
(17)
vapor at low pressure, Nf is about 1017 per cubic centimeter, ω0 is about 3 × 1015 per second, and g is about 1011 per second. These data show that, when the frequency of the radiation is not very near ω0, ωg is very small in comparison with the denominator and the absorption is practically zero. As ω approaches ω0, κ increases rapidly to a maximum at a radiational frequency very near ω0 and then falls at frequencies greater than ω0. When the absorption is relatively weak, the absorption maximum is directly proportional to the oscillator strength f. In terms of the molar extinction coefficient of Eq. (4), it can be shown that this direct relation holds, as seen in Eq. (18). The integration in Eq. (18) is carried out f = 4.319 × 10−9 d ν (18) over the whole absorption spectrum. The integral can be evaluated from the area under the curve of plotted as a function of wave number ν cm−1 = ν (s−1)/c = 1/λ. The width of the absorption band for an atom is determined by the value of the damping coefficient g; the greater the damping, the greater is the spectral region over which absorption extends. The general behavior of the refractive index through the absorption band is illustrated by the dotted portions of Fig. 2. The presence of the damping term ω2g2 in the denominator of Eq. (17) prevents
(15) 2.0
radio waves
1.0
n
So long as the absorption remains negligible, these equations correctly describe the increase in refractive index as the frequency of the radiation begins to approach the absorption band determined by ωi or λi. They fail when absorption becomes appreciable, since they predict infinitely large values of the refractive index when ω equals ωi, whereas the refractive index remains finite throughout an absorption band.
far visible ultraviolet
0
x-rays
far- infrared near- infrared
nearultraviolet λ
Fig. 2. Complete dispersion curve through the electromagnetic spectrum for a substance. (After F. A. Jenkins and H. E. White, Fundamentals of Optics, 4th ed., McGraw-Hill, 1976)
29
30
Abstract algebra the refractive index from becoming infinite when ω = ω0. Its value increases to a maximum for a radiation frequency less than ω0, then falls with increasing frequency in the center of the band (anomalous dispersion) and increases from a relatively low value on the high-frequency side of the band. Figure 2 shows schematically how the dispersion curve is determined by the absorption bands throughout the whole electromagnetic spectrum. The dotted portions of the curve correspond to absorption bands, each associated with a distinct type of electrical oscillator. The oscillators excited by x-rays are tightly bound inner elections; those excited by ultraviolet radiation are more loosely bound outer electrons which control the dispersion in the near-ultraviolet and visible regions, whereas those excited by the longer wavelengths are atoms or groups of atoms. It will be observed in Fig. 2 that in regions of anomalous dispersion the refractive index of a substance may assume a value less than unity; the velocity of light in the medium is then greater than in vacuum. The velocity involved here is that with which the phase of the electromagnetic wave of a single frequency ω advances, for example, the velocity with which the crest of the wave advances through the medium. The theory of wave motion, however, shows that a signal propagated by electromagnetic radiation is carried by a group of waves of slightly different frequency, moving with a group velocity which, in a material medium, is always less than the velocity of light in vacuum. The existence of a refractive index less than unity in a material medium is therefore not in contradiction with the theory of relativity. In quantum theory, absorption is associated not with the steady oscillation of a charge in an orbit but with transitions from one quantized state to another. The treatment of dispersion according to quantum theory is essentially similar to that outlined, with the difference that the natural frequencies ν 0 are now identified with the frequencies of radiation which the atom can absorb in undergoing quantum transitions. These transition frequencies are regarded as virtual classical oscillators, which react to radiation precisely as do the oscillators of classical electromagnetic theory. Selective reflection. Nonmetallic substances which show very strong selective absorption also strongly reflect radiation of wavelengths near the absorption bands, although the maximum of reflection is not, in general, at the same wavelength as the maximum absorption. The infrared rays selectively reflected by ionic crystals are frequently referred to as reststrahlen, or residual rays. For additional information on selective reflection. See REFLECTION OF ELECTROMAGNETIC RADIATION. William West Bibliography. M. Born and E. Wolf, Principles of Optics, 7th ed., 1999; R. W. Ditchburn, Light, 3d ed., 1977, reprint 1991; I. S. Grant and W. R. Phillips, Electromagnetism, 2d ed., 1991; B. D. Guenther, Modern Optics, 1990; E. Hecht, Optics, 4th ed., 2001; F. A. Jenkins and H. E. White, Fundamen-
tals of Optics, 4th ed., 1976; E. E. Kriezis, D. P. Chrissoulidis, and A. C. Papagiannakis (eds.), Electromagnetics and Optics, 1992; Optical Society of America, Handbook of Optics, 2d ed., 4 vols., 1995, 2001; A. Sommerfeld, Lectures of Theoretical Physics, vol. 4, 1954.
Abstract algebra The study of systems consisting of arbitrary sets of elements of unspecified type, together with certain operations satisfying prescribed lists of axioms. Abstract algebra has been developed since the mid1920s and has become a basic idiom of contemporary mathematics. In contrast to the earlier algebra, which was highly computational and was confined to the study of specific systems generally based on real and complex numbers, abstract algebra is conceptual and axiomatic. A combination of the theories of abstract algebra with the computational speed and complexity of computers has led to significant applications in such areas as information theory and algebraic geometry. Insight into the difference between the older and the abstract approaches can be obtained by comparing the older matrix theory with the more abstract linear algebra. Both deal with roughly the same area of mathematics, the former from a direct approach which stresses calculations with matrices, and the latter from an axiomatic and geometric viewpoint which treats vector spaces and linear transformations as the basic notions, and matrices as secondary to these. See LINEAR ALGEBRA. Algebraic structures. Abstract algebra deals with a number of important algebraic structures, such as groups, rings, and fields. An algebraic structure consists of a set S of elements of unspecified nature endowed with a number of finitary operations on S. If r is a positive integer, an r-ary operation associates with every r-tuple (a1, a2,. . . , ar) of elements ai in S a unique element ω(a1, a2,. . . , ar) in S. It is convenient to consider also nullary operations, that is, the selection of constants from S. Binary operations. Many of the operations considered in abstract algebra are binary (2-ary) operations, in which case the symbol for the operation is usually written between the elements, or juxtaposition is used to indicate the operation, so that ω(a, b) is typically represented by a ∗ b, a + b, or ab, depending on the context. Algebraic structures are classified by properties that the operations of the structure may or may not satisfy. Among the most common properties are: 1. Associativity: A binary operation ∗ of an algebraic structure S is associative if a ∗ (b ∗ c) = (a ∗ b) ∗ c for all elements a, b, c of S. 2. Commutativity: A binary operation ∗ of an algebraic structure S is commutative if a ∗ b = b ∗ a for all elements a, b of S. 3. Existence of an identity: An element e of an algebraic structure S is an identity for the binary
Abstract algebra operation ∗ if a ∗ e = e ∗ a = a for all elements a of S. 4. Existence of inverses: Assuming there is an element e of an algebraic structure S that is an identity for the binary operation ∗, an element a of S has an inverse if there is an element a −1 of S such that a ∗ a −1 = a −1 ∗ a = e. 5. Distributivity: A binary operation ∗ distributes over a binary operation + of an algebraic structure S if a ∗ (b + c) = (a ∗ b) + (a ∗ c) for all elements a, b, c of S. The property of distributivity thus provides a link between two operations. Groups. If the structure S is a group, a single binary operation is assumed to satisfy several simple conditions called the group axioms: the associativity of the operation, the existence of an identity element for the operation, and the existence of inverses for the operation. In this case, ab is usually written for ω(a,b), and the operation is called multiplication; or a +b is written, and the operation is called addition, if the operation ω is commutative. Alternatively, a group may be described as a set with three operations rather than only one. In addition to the associative binary operation mentioned above, it is possible to consider a unary operation providing the inverse of an element, and a nullary operation providing the identity element of the group. Examples of groups include the rigid motions (symmetries) of a square with operation given by function composition, the integers under addition, and the nonzero real numbers under multiplication. Rings and fields. In the case of a ring R, there are two binary operations, ab and a + b, which are assumed to satisfy a set of conditions known as the ring axioms, ensuring the associativity of multiplication and addition, the commutativity of addition, the distributivity of multiplication over addition, and the existence of an additive identity and additive inverses. Examples of rings include the integers, the real numbers, and the continuous real-valued functions of a real variable. See RING THEORY. If axioms asserting the commutativity of multiplication and the existence of a multiplicative identity and multiplicative inverses are adjoined to the ring axioms, the ring is also a field. Examples of fields include the real numbers, the complex numbers, the finite set of integers modulo a prime number, and the rational polynomials with complex coefficients. See FIELD THEORY (MATHEMATICS). Modules. In addition to the intrinsic study of algebraic structures, algebraists are interested in the study of the action of one algebraic structure on another. An important instance is the theory of modules and its special case of vector spaces. A left module over a ring R is defined to be a commutative group M on which R acts on the left in the sense that, given a pair of elements (a,x) where a is in R and x is in M, then the action determines a unique element ax in M. Moreover, the module product ax is assumed to satisfy the module axioms a (x + y) = ax + ay, (a + b) x = ax + bx, and (ab)x = a (bx), where a and b are arbitrary elements of R, and x and y are arbitrary elements of M. Consideration of the modules
admitted by a ring can provide significant information about the structure of the ring. Specializing from an arbitrary ring to a field yields the special type of module called a vector space. Comparison of similar structures. Along with the study of particular algebraic structures, algebraists are interested in comparing structures satisfying similar axioms or sharing certain properties. Those structures most intimately related to a particular structure are its substructures and quotient structures. Substructures. A substructure of an algebraic structure S is a subset T that is closed with respect to the operations of S; that is, if ω is any r-ary operation of S and a1, a2, . . . , ar are elements of the subset T of S, then ω(a1, a2, . . . , ar) is again an element of T. In this case, T with operations defined by the operations of S restricted to T becomes an algebraic structure itself. Equivalently S may be referred to as an extension of T. As an example, if S is taken as the commutative group of integers under addition, the subset T of multiples of any particular integer forms a subgroup of S. Quotient structures. A quotient structure is defined on an algebraic structure S by partitioning the elements of S into a family Q of (possibly infinitely many) disjoint congruence classes Qα. As a result, (1) each element of S is a member of one and only one of the subsets Qα of S, and (2) for any r-ary operation ω of S, no matter what elements a1, a2, . . . , ar of Qα1, Qα2, . . . , Qαr respectively are chosen, the subset Qβ of S containing ω(a1, a2, . . . , ar) is the same. Then the collection of congruence classes S/Q = {Qα} can be made into an algebraic structure with operations determined by the action of the operations of S on the elements of the congruence classes, called a quotient structure of S. For example, an r-ary operation ω would act on an r-tuple (Qα1, Qα2, . . . , Qαr) by assigning to it the congruence class Qβ if it contains the element ω(a1, a2, . . . , ar) for elements a1, a2, . . . , ar of Qα1, Qα2, . . . , Qαr, respectively. An example of a quotient structure is one in which S is the ring of integers under addition and multiplication, and the family of congruence classes Q consists of the subset E of even integers and the subset O of odd integers. Then the quotient structure S/Q contains two elements E and O; the operations are given by E + E = E = O + O, E + O = O = O + E, E ∗ E = E = E ∗ O = O ∗ E, and O ∗ O = O. These operations symbolically represent the familiar facts that the sum of two even or two odd integers is even, the sum of an even integer and an odd integer is odd, the product of an even integer with either an even or an odd integer is even, and the product of two odd integers is odd. Comparisons by means of functions. Other comparisons between similar algebraic structures are made by means of functions between the structures that preserve the action of the corresponding operations. Two principal ways of organizing this study are universal algebra and category theory. 1. Universal algebra. A considerable part of the study of algebraic structures can be developed in a unified way that does not specify the particular
31
32
Abstract data type structure, although deeper results do require specialization to the particular systems whose richness is to a large extent accounted for by their applicability to other areas of mathematics and to other disciplines. The general study of algebraic structures is called universal algebra. A basic notion here is that of a homomorphism of one algebraic structure S to a second one S . Here there is a one-to-one correspondence ω ↔ ω between the sets of operations of S and of S such that ω and ω are both r-ary for the same r = 0, 1, 2, 3, . . . . Then a homomorphism f from S to S is a function mapping S to S such that ω(a1, . . . , ar) = ω(f(a1),. . . , f (ar)) for all ai in S and all the corresponding operations ω and ω. A substantial part of the theory of homomorphisms of groups and rings is valid in this general setting. Another basic notion of universal algebra is that of a variety of algebraic structures, which is defined to be the class of structures satisfying a given set of identities defined by means of the given compositions (that is, the associative or the commutative laws). Examples are the variety of groups and the variety of rings. 2. Category theory. A second approach to the comparison of related structures is that of category theory. The notion of a category is applicable throughout mathematics, for example, in set theory and topology. The notion is formulated in order to place on an equal footing a class of mathematical objects and the basic mappings between these; in contrast is the more elemental focus of universal algebra. A category is defined to be a class of objects (generally mathematical structures) together with a class of morphisms (generally mappings of the structures) subject to certain axioms. An example of a category is the class of all groups as objects and the class of all homomorphisms of pairs of groups as the class of morphisms. Another example is the class of all sets as objects and of mappings of one set into another as the class of morphisms. As an example of the ever-increasing levels of abstraction available and useful in algebra, once categories have been defined, they can themselves be compared by means of functors, which act as homomorphisms of categories. If C and C are categories, a functor from C to C is defined as a mapping F of the objects of C into the objects of C , together with a mapping φ of the morphisms of C into the morphisms of C , such that φ preserves the compositions and identity morphisms. At this point, functors can be compared via natural transformations, which thus act as homomorphisms of functors. The concepts of category, functor, and natural transformation appear in many areas of mathematics. Many of the homomorphisms that arise in the study of algebraic structures prove to be instances of natural transformations. The frequency of the occurrences of the naturalness of the homomorphisms has implications for any adequate philosophy of mathematics. Like set theory, category theory can provide a foundation for much of mathematics. Joel K. Haack Bibliography. M. Artin, Algebra, 1995; J. R. Durbin, Modern Algebra: An Introduction, 4th ed., 1999; J. Gallian, Contemporary Abstract Algebra, 2d ed.,
1990; T. W. Hungerford, Algebra, rev. ed., 1980, reprint 1993; S. MacLane and G. Birkhoff, Algebra, 3d ed., 1987; C. C. Pinter, A Book of Abstract Algebra, 2d ed., 1990.
Abstract data type A mathematical entity consisting of a set of values (the carrier set) and a collection of operations that manipulate them. For example, the Integer abstract data type consists of a carrier set containing the positive and negative whole numbers and 0, and a collection of operations manipulating these values, such as addition, subtraction, multiplication, equality comparison, and order comparison. See ABSTRACT ALGEBRA. Abstraction. To abstract is to ignore some details of a thing in favor of others. Abstraction is important in problem solving because it allows problem solvers to focus on essential details while ignoring the inessential, thus simplifying the problem and bringing to attention those aspects of the problem involved in its solution. Abstract data types are important in computer science because they provide a clear and precise way to specify what data a program must manipulate, and how the program must manipulate its data, without regard to details about how data are represented or how operations are implemented. Once an abstract data type is understood and documented, it serves as a specification that programmers can use to guide their choice of data representation and operation implementation, and as a standard for ensuring program correctness. A realization of an abstract data type that provides representations of the values of its carrier set and algorithms for its operations is called a data type. Programming languages typically provide several built-in data types, and usually also facilities for programmers to create others. Most programming languages provide a data type realizing the Integer abstract data type, for example. The carrier set of the Integer abstract data type is a collection of whole numbers, so these numbers must be represented in some way. Programs typically use a string of bits of fixed size (often 32 bits) to represent Integer values in base two, with one bit used to represent the sign of the number. Algorithms that manipulate these strings of bits implement the operations of the abstract data type. See ALGORITHM; DATA STRUCTURE; NUMERICAL REPRESENTATION (COMPUTERS); PROGRAMMING LANGUAGES. Realizations of abstract data types are rarely perfect. Representations are always finite, while carrier sets of abstract data types are often infinite. Many individual values of some carrier sets (such as real numbers) cannot be precisely represented on digital computers. Nevertheless, abstract data types provide the standard against which the data types realized in programs are judged. Example. Data processing problems are rarely so simple that data abstraction is not helpful in their solution. Suppose a program must process data in its
Acantharea order of arrival but that data may arrive before its predecessors have been completely processed. Data values must wait their turn in storage in a way that preserves their arrival order. Before considering the structure needed to store the data and the algorithms required for creating and maintaining the structure, it is useful to specify an abstract data type for this problem. What is needed is a well-known abstract data type called a Queue. The carrier set of the Queue abstract data type is the set of all queues, that is, the set of all lists of data values whose elements enter at the rear and are removed from the front. Although many operations might be included in this abstract data type, the following four suffice: 1. The enter operation takes a queue and a data value as parameters and returns the input queue with the new data value added to its rear. 2. The remove operation takes a queue as its parameter and returns the input queue with the data value at its front removed. 3. The front operation takes a queue as its parameter and returns the data value at the front of the queue (the queue is unaltered). 4. The size operation takes a queue as its parameter and returns a whole number greater than or equal to 0 reporting the number of data values in the queue. These informal operation descriptions should be augmented by specifications that precisely constrain them: 1. For any queue q, if size (q) = 0 then remove (q) is invalid. That the operation is invalid is reported as an error. 2. For any queue q, if size (q) > 0 then size (q) = size (remove (q)) + 1. 3. For any queue q and data value v, size (q) = size (enter (q,v)) − 1. 4. For any queue q and data value v, let q = enter (q,v). Then after an arbitrary sequence of enter and remove operations applied to q containing exactly size (q) remove operations, and yielding q, v = front (q). The first of these specifications states that nothing can be removed from an empty queue. The second and third state how queues grow and shrink. The final specification captures the crucial property that all data values eventually leave a queue in the order that they enter it. Some reasoning will show that these four specifications completely determine the behavior of queues. Usefulness. Such specifications of abstract data types provide the basis for their realization in programs. Programmers know which data values need to be represented, which operations need to be implemented, and which constraints must be satisfied. Careful study of program code and the appropriate selection of tests help to ensure that the programs
are correct. Finally, specifications of abstract data types can be used to investigate and demonstrate the properties of abstract data types themselves, leading to better understanding of programs and ultimately higher-quality software. See COMPUTER PROGRAMMING; SOFTWARE ENGINEERING; SHRINK FIT. Relation to object-oriented paradigm. A major trend in computer science is the object-oriented paradigm, an approach to program design and implementation using collections of interacting entities called objects. Objects incorporate both data and operations. In this way they mimic things in the real world, which have properties (data) and behaviors (operations). Objects that hold the same kind of data and perform the same operations form a class. Abstract data values are separated from abstract data type operations. If the values in the carrier set of an abstract data type can be reconceptualized to include not only data values but also abstract data type operations, then the elements of the carrier set become entities that incorporate both data and operations, like objects, and the carrier set itself is very much like a class. The object-oriented paradigm can thus be seen as an outgrowth of the use of abstract data types. See OBJECT-ORIENTED PROGRAMMING. Christopher Fox Bibliography. J. Guttag, Abstract data types and the development of data structures, Commun. ACM, 20(6):396–404, June 1977; P. Thomas, H. Robinson, and J. Emms, Abstract Data Types, Oxford University Press, Oxford, 1988.
Acantharea A class of the phylum Radiozoa in the Actinopoda. The kingdom Protozoa contains 18 phyla. One of the parvkingdoms (a hierarchical classification between kingdom and superphylum that is controversial and not officially recognized as such) is the Actinopoda (originally a class) containing two phyla: Heliozoa and Radiozoa. Within the Radiozoa is the class Acantharea. These marine protozoans possess a nonliving, organic capsular wall surrounding a central mass of cytoplasm. The intracapsular cytoplasm is connected to the extracellular cytoplasm by fine cytoplasmic strands passing through pores in the capsular wall. When viewed with the electron microscope, the capsular wall in some species appears to be made of many layers, each composed of a fine fibrillar network. The skeletons are made of celestite (strontium sulfate) instead of silica. The basic structural elements are 20 rods that pass through the capsule to the center in regular arrangements (polar and equatorial; see illustration). An equatorial rod forms a 90◦ angle with a polar rod, and other groups are arranged with similar exactness. This type of cytoskeleton may be modified by the addition of a latticework, composed of plates, each fused with a skeletal rod. Some genera show a double latticework, concentric with the central capsule. The skeletons do not contribute substantially to the paleontological record in marine sediments since celestite is dissolved by seawater.
33
34
Acanthocephala
central capsule membrane axial rod
central capsule myonemes (a)
(b)
(c)
(d) Representatives of the Acantharea; Acanthometra pellucida, (a) showing axial rods extending through the central capsule, and myonemes (myophrisks) extending from the rods to the sheath; (b) contracted myophrisks with expanded sheath; and (c) relaxed myophrisks with contracted sheath. (d) Dorotaspis heteropora, with perforated shell and rods. (After R. P. Hall, Protozoology, Prentice-Hall, 1953)
Dissolution is pronounced below 617 ft (188 m) depth. While the protozoan is still alive, however, the cytoplasm appears to protect the skeleton from dissolution. Pseudopodia are more or less permanent. A pseudopodium either is composed of an axial rod containing microtubules and a cytoplasmic envelope (forming an axopodium, or small foot; long slender pseudopodia are found in certain protozoa) or is finebranched (forming a reticulopodium). Mitochondria are always present and contain flattened cristae. No cilia are present in the tropic phase. The central endoplasm contains numerous nuclei and other organelles. When present, zooxanthellae (algae found within the cytoplasm) are of at least two kinds: Dinophyceae containing trichocysts in the cytoplasm; and a group of algae characterized by numerous discoid plastids, each of which contains an interlamellar pyrenoid. The systematics of the second group is still uncertain. Myonemes (myophrisks) are significant components of the hydrostatic apparatus and regulate the buoyancy of the Acantharea by expanding or contracting portions of the extracapsular sheath. When observed with the electron microscope, they are composed of microfilaments arranged in a cone around the spine. Microfilaments contain actin, a contractile protein also found in metazoan muscle. Their contraction extends the sheath (illus. b); their relaxation flattens it against the rods (illus. c). Although essentially pelagic, Acantharea may move vertically with the help of their hydrostatic ap-
paratus. Symbiont-bearing species (species that contain zooxanthellae) have been reported from depths as great as 13,120 ft (4000 m). Little is known about the ecology or distribution of the Acantharea. The Gulf Stream is rich with these protozoans in the spring and summer but poor during other seasons. When compared to other plankton, Acantharea have a low average percentage of abundance even at peak periods (about 6–7%). Their approximate average abundance in the North Atlantic has been reported as 1–5%. See GULF STREAM. Information about life cycles is very limited. For a few species, flagellate stages, syngamy, and zygotes have been reported. Certain stages of fission have been reported for some primitive types. The taxonomy of the Acantharea is still under review. See ACTINOPODEA; CELESTITE; PROTOZOA; SARCODINA; SARCOMASTIGOPHORA; STRONO. Roger Anderson; John P. Harley TIUM. Bibliography. T. Cavalier-Smith, Kingdom Protozoa and its 18 phyla, Microbiol. Rev., 57(4):953–994, 1993; M. Grell, in H.-E. Gruner (ed.), Lehrbuch der Speziellen Zoologie, 4th ed., Gustav Fischer, Stuttgart, 1980; K. Hausmann and N. Hulsmann, Protozoology, Georg Thieme Verlag, New York, 1996.
Acanthocephala A distinct phylum of helminths, the adults of which are parasitic in the alimentary canal of vertebrates. They are commonly known as the spiny-headed worms. The phylum comprises the orders Archiacanthocephala, Palaeacanthocephala, and Eoacanthocephala. Over 500 species have been described from all classes of vertebrates, although more species occur in fish than in birds and mammals and only a relatively few species are found in amphibians and reptiles. The geographical distribution of acanthocephalans is worldwide, but genera and species do not have a uniform distribution because some species are confined to limited geographic areas. Host specificity is well established in some species, whereas others exhibit a wide range of host tolerance. The same species never occurs normally, as an adult, in cold-blooded and warm-blooded definitive hosts. More species occur in fish than any other vertebrate; however, Acanthocephala have not been reported from elasmobranch fish. The fact that larval development occurs in arthropods gives support to the postulation that the ancestors of Acanthocephala were parasites of primitive arthropods during or before the Cambrian Period and became parasites of vertebrates as this group arose and utilized arthropods for food. See ARCHIACANTHOCEPHALA; EOACANTHOCEPHALA; PALAEACANTHOCEPHALA. Morphology Adults of various species show great diversity in size, ranging in length from 0.04 in. (1 mm) in some species found in fish to over 16 in. (400 mm) in some mammalian species (Fig. 1). Most of the larger
Acanthocephala
(a)
(b)
(d)
(e)
(c)
(f)
Fig. 1. Various acanthocephalan examples, adult forms drawn to same scale. (a) Moniliformis dubius, 4–12 in. (10–30 cm) long. (b) Proboscis of same. (c) Macracanthorhynchus hirudinaceus, 10–24 in. (25–60 cm) long. (d) Proboscis of same. (e) Oncicola canis, 0.5–0.8 in. (1–2 cm) long. (f) Proboscis of same. (After A. C. Chandler and C. P. Read, Introduction to Parasitology, 10th ed., Wiley, 1961)
species are from mammals; however, some mammalian forms are relatively small (Corynosoma) even though they parasitize seals and whales. When observed undisturbed in the intestine, the elongate body is distinctly flattened, but on removal to water or saline and after preservation the body becomes cylindrical. Living worms are translucent or milky white, although they may take on a distinctive color from the intestinal contents. The body of both sexes has three divisions: the proboscis armed with hooks, spines, or both; an unspined neck; and the posterior trunk. Proboscis. The proboscis is the primary organ for attachment of Acanthocephala to the intestinal wall of the host. In some species the proboscis is parallel with the main axis of the body, but frequently it is inclined ventrally. The proboscis may be globular, or elongate and cylindrical in shape and is invariably armed with sharp pointed hooks, spines, or both (Fig. 1). These structures vary in shape and number and are usually arranged radially or in spiral rows. Electron microscopy studies have revealed that proboscis hooks consist of a central core of cytoplasm covered by hard nonliving material. In species with an elongate proboscis the hooks seem to be in longitudinal rows with a quincuncial arrangement. In most species the proboscis is capable of introversion into a saclike structure, the proboscis receptacle. When in this position, the tips of the hooks are all directed anteriorly and are on the inside of the introverted proboscis. As the proboscis is everted and extended, these hooks engage the mucosa of the host’s intestine, with the tips turning posteriorly as they reach their functional position on the exterior of the proboscis. The recurved points of the hooks anchor the worm firmly to the intestinal wall (Fig. 2). The proboscis receptacle and neck can be retracted into the body cavity but without inversion.
Within the anterior end of the trunk are found the proboscis receptacle; musculature for retracting the proboscis, receptacle, and neck; and the lemnisci (Fig. 3). The proboscis receptacle is attached to the inner surface of the proboscis. Muscles known as inverters of the proboscis are attached to the anterior tip of this structure and pass through the proboscis and receptacle, emerging from the posterior of the receptacle, and continue some distance posteriorly, where they attach to the trunk wall. Usually one inverter is attached dorsally and one ventrally. Contraction of the inverter fibers introverts the proboscis into the receptacle, while further contraction of these fibers pulls the receptacle deep into the body cavity. Body wall. Electron microscopy studies show that the body wall is covered by a thin mucopolysaccharide known as the epicuticle secreted onto the surface of the worm. Beneath the epicuticle is the thin cuticle with an outer membrane and a homogeneous layer perforated by numerous pores of canals which pass through the striped layer beneath. Many of the canals are filled with an electron-dense material. Below the striped layer is the felt layer composed of many fibrous strands extending in various directions. Mitochondria and numerous vesicles are present in the felt layer. Some vesicles seem to be connected to the canals of the striped layer. The radial layer beneath the felt layer contains fewer fibrous strands, large thin-walled lacunar channels with mitochondria arranged around the channels. Lipids and glycogen are present in the radial layer. Circular and longitudinal muscle layers are found beneath the radial layer separated by plasma membranes. Some investigators consider all body-wall layers beneath the stripped layer as hypodermis.
host intestine
(c)
(b)
(a)
(d)
(e)
Fig. 2. Diagrams showing mechanics of proboscis attachment of Neoechinorhynchus emydis in intestinal wall of a turtle. (a) Proboscis fully introverted. (b) Most basal hooks in contact with host tissue. (c–e) Progressive stages in extroversion of the proboscis. (After H. J. Van Cleave, Experimental Parasitology, vol. 1, 1952)
proboscis
brain
proboscis receptacle
lemnisci
retractor muscles
Fig. 3. Anterior end of Moniliformis dubius.
35
36
Acanthocephala
remnant of the ligament
fragmented epidermal nuclei epidermis
circular muscle layer
cuticle pseudocoel
main lacunar channels
longitudinal muscle layer
lemnisci (a) epidermis
giant epidermal nucleus
dorsal ligament sac cuticle
circular muscle layer pseudocoel longitudinal muscle layer
Reproductive system. The reproductive organs of the male consist of a pair of testes and specialized cells, the cement glands. In most species there is a saclike structure behind the cement glands, Saefftigen’s pouch, through which run the sperm ducts and the ducts from the cement glands (Fig. 6a). The products of the testes and cement glands are discharged through a penis, which is surrounded by a posteriorly located bell-shaped copulatory bursa which is usually introverted into the posterior extremity of the trunk. At copulation the bursa is extended and applied to the posterior extremity of the female, where it is held firmly in place by the secretions of the cement glands. This material hardens, forming a copulatory cap on the female which is an internal cast of the bursa of the male, and remains attached to the posterior extremity of the female for some time following copulation.
ventral ligament sac (b)
main lacunar channels
Fig. 4. Cross sections showing body plans of acanthocephalans. (a) Palaeacanthocephalan. (b) Archiacanthocephalan. (After L. H. Hyman, The Invertebrates, vol. 3, McGraw-Hill, 1951)
testes
vasa deferentia Fig. 5. Lacunar system of Moniliformis, dorsal view, showing regular circular branches. (After L. H. Hyman, The Invertebrates, vol. 3, McGraw-Hill, 1951)
Giant nuclei. The wall of the trunk has an external cuticula beneath which lies the thick syncytial hypodermis, or subcuticula. The hypodermis contains a relatively small number of giant nuclei which may be round, oval, ameboid, or elongate with a number of lateral branches (Fig. 4). Lacunae. Within the hypodermis is a series of intercommunicating spaces, the lacunar system. Usually the lacunar system is composed of two longitudinal vessels either dorsal and ventral or ventral and lateral. In some species only the dorsal vessel is present. The longitudinal vessels are connected by a series of smaller vessels, the lacunae, which ramify throughout the body. In species exhibiting pseudosegmentation (Moniliformis), the regularly spaced lateral lacunae and the body musculature divide the wall into transverse folds (Fig. 5). These folds have no effect on the arrangement of internal organs. Pseudocoele. The body cavity, or pseudocoele, contains all the internal organs, the most conspicuous of which are the reproductive organs enclosed in axial connective-tissue ligament sacs. These ligament sacs are hollow tubes which extend most of the length of the cavity of the trunk. They are single in both males and females of the Palaeacanthocephala, but are divided into dorsal and ventral sacs which communicate anteriorly in females of the other orders. There is no vestige of a digestive system in any stage of the life cycle.
testes vasa deferentia
(a)
cement glands
Saefftigen's pouch dorsal ligament sac ventral ligament sac uterine bell
copulatory bursa
uterine pouches sorting apparatus uterus sphincter
genital orifice
vagina (b) Fig. 6. Reproductive system of Moniliformis dubius. (a) Posterior end of male. (b) Posterior end of female. (After A. C. Chandler and C. P. Read, Introduction to Parasitology, 10th ed., Wiley, 1961)
Acanthocephala Female Acanthocephala are unique in that the ovary exists as a distinct organ only in the very early stages of development and later breaks up to form free-floating egg balls. The eggs are fertilized as they are released from the egg balls and are retained within the ligament sacs until embryonation is complete. The genital orifice is posterior. A vagina provided with a strong sphincter extends anteriorly from the orifice and a saccular uterus is anterior to the vagina. The anterior end of the uterus is surrounded by a series of guard cells, the selective apparatus. From this extends a funnellike structure, the uterine bell, the broad anterior end of which opens into the body cavity or one of the ligament sacs holding the developing eggs (Fig. 6b). During embryonation in the body cavity the eggs acquire a series of membranes and a hard outer shell. Eggs which have not completed embryonation are passed back into the body cavity through special openings in the selective apparatus, whereas eggs which are fully mature are passed into the saccular uterus and eliminated through the vaginal sphincter and genital orifice into the intestinal contents of the host. Lemnisci. The trunk and presoma are demarcated by an infolding of the cuticula and the presence of two elongate contractile structures, the lemnisci, which arise from the hypodermis and extend posteriorly into the body cavity. The lemnisci are provided with a definite number of nuclei which migrate from the hypodermis into the lemnisci as they are formed during larval development. The function of the lemnisci is unknown. One explanation offered for their function is that they act as reservoirs for the fluid of the lacunar system when the presoma and proboscis are invaginated. Nervous system. The nervous system is composed of a chief ganglion or brain located within the proboscis receptacle. Most of the nerve trunks are inconspicuous but two of them, the retinacula, associated with muscle fibers, pass through the wall of the proboscis receptacle to innervate the trunk wall. Excretory system. In members of the Archiacanthocephala modified protonephridial organs are found closely adherent to the reproductive system, but in most species specialized excretory organs are completely lacking. The protonephridial organs consist of a mass of flame bulbs attached to a common stem. The canals unite to form a single canal which joins the sperm duct in the male and the uterus in the female. Embryology The Acanthocephala are obligatory parasites throughout their entire life cycle; no known member of this phylum exists as a free-living organism. Acanthor. The eggs passed in the feces of the host contain a mature embryo, the acanthor, which is surrounded by three membranes (Fig. 7a). The life cycle always involves an intermediate host, which is usually an arthropod. Among these are small crustaceans for parasites of aquatic vertebrates, and
grubs, roaches, and grasshoppers for those which parasitize terrestrial vertebrates. The eggs hatch after being ingested by the appropriate arthropod and release the spindle-shaped acanthor, which has a spiny body and is armed with retractile bladelike rostellar hooks used in the penetration of the intestinal wall of the intermediate host (Fig. 7b). The central portion of the body of the acanthor contains a column of dense undifferentiated nuclei, the entoblast, from which develop the structures of the ensuing stages. After penetration of the wall of the digestive tract, the acanthor comes to lie beneath the delimiting membrane of the gut wall, becomes quiescent, and begins to grow (Fig. 7c). It soon drops free into the host’s hemocoele, where it undergoes a gradual transformation into the various larval stages (Fig. 7d and e). Acanthella. The term acanthella is applied to the series of stages in which the rudiments of the reproductive organs, lemnisci, proboscis, and proboscis receptacle are formed (Fig. 7f). As development progresses, the acanthella becomes elongate, and is surrounded by a hyaline cyst produced by the larva. Cystacanth. When the proboscis and the proboscis receptacle of the acanthella develop to the point of retractility, the larva becomes infective and is known as the cystacanth (Fig. 7g and h). The mature or infective cystacanth lies in the hemocoele of the intermediate host with the proboscis and receptacle retracted into the body cavity. The body cavity of the cystacanth contains all of the structures of the adult worm in an immature form. In the least complicated life cycle the cystacanth remains in the hemocoele of the arthropod until it is ingested by a suitable vertebrate host, in which it excysts and attaches to the intestinal wall by the proboscis and grows to sexual maturity. In some forms the life cycle may be prolonged and complicated by the introduction of one or more transport hosts, usually a vertebrate, in which the cystacanth becomes established as a visceral cyst awaiting ingestion by a suitable vertebrate host in which sexual maturity can be attained. Cell constancy. The somatic tissues of Acanthocephala are made up of a restricted number of nuclei which is more or less fixed for each tissue. This condition is termed cell constancy or more correctly eutely or nuclear constancy, because the majority of acanthocephalan tissues are syncytial in nature. The nuclei of the early embryonic cells which are destined to become syncytial in nature are large, regular, and spheroidal or elliptical in shape. During the transformation of the acanthella into the cystacanth the nuclei begin to assume the shape and form of the hypodermal nuclei seen in the adult by passing through an integrated series of ovoid, rosette, ameboid, elongate with lateral branches, and arboroid or dendritic forms. These diverse expressions of nuclear shape and modifications are instrumental in maintaining a favorable ratio between nuclear surface and surrounding cytoplasm. At the conclusion of the cleavage process the embryo, or acanthor, contains all of
37
38
Acanthocephala
0.05 mm
outer shell
rostellar hooks
rostellar hooks
retractor muscles cuticular spines
inner shell
central nuclear mass
central nuclear mass
0.02 mm
(c) proboscis
(b)
nuclei of lemniscal ring
(a) 0.05 mm
proboscis receptacle lemnisci
proboscis proboscis receptacle
brain
brain
retractor muscles
dermal layer giant nuclei (d)
nuclei of apical ring
giant nuclei 0.2 mm
inverters
uterine bell
nuclei of apical ring
(e) 0.5 mm
proboscis receptacle
0.2 mm
brain cyst
lemnisci
retractor muscles
giant nuclei
genital ligament
0.5 mm
testes
(f)
(g)
cement glands
(h)
Fig. 7. Life history of Moniliformis dubius. (a) Mature egg containing acanthor. (b) Acanthor in process of escaping from egg shells and membranes. (c) Acanthor dissected from gut wall tissue of Periplaneta americana 5 days after infection. (d) Median sagittal section of larva removed from the body cavity of P. americana 29 days after infection. (e) Larva removed from body cavity of P. americana 42 days after infection. (f) Acanthella dissected from its enveloping sheath. (g) Encysted cystacanth from body cavity of P. americana with proboscis invaginated. (h) Cystacanth freed from its cyst and with proboscis evaginated.
Acanthodii the cellular elements of the adult. No mitotic divisions take place, although in the rearrangement of nuclei during larval development certain nuclei may undergo amitotic divisions. The extent of the metamorphoses of the embryonic cells and nuclei varies greatly in the different groups and genera. In the more highly specialized genera, the nuclei of the embryo lose all of their original appearance during larval development. In contrast, the nuclei of the adults of the less specialized genera retain the distinctive embryonic nuclei with little change. Varying changes, intermediate between these two extremes, are observed in other groups. See CELL CONSTANCY. Physiology Since acanthocephalans lack a gut, they undoubtedly obtain nutrients through the metasomal surface from material in the host’s intestinal lumen. The role of the pores in this process is not really known, but it is likely that absorption of nutrients is facilitated by the increased surface area resulting from the plasma membrane lining the pores. It has been suggested that digestion may take place in the region of the presoma, which is in continuous contact with the host tissue. In this area, digestive enzymes are believed to be membrane-bound, allowing products of digestion to be absorbed at the host-parasite interface. Acanthocephalans contain large quantities of glycogen and in general have a pronounced carbohydrate metabolism. Glycogen is localized in the proboscis, body wall, and muscles. They appear to utilize the Embden-Meyerhoff scheme of phosphorylating glycolysis. Chief waste products of carbohydrate metabolism in Moniliformis dubius are succinic, lactic, formic, and acetic acids, along with large quantities of ethanol. It has been suggested that the end products of carbohydrate metabolism excreted by acanthocephalans may be metabolized by the host or may be effective in preventing other helminths from establishing in that area of the host’s intestine. See CARBOHYDRATE METABOLISM; GLYCOGEN. A significant Krebs cycle has not been demonstrated in acanthocephalans, but evidence indicates that it may be partially operative in some species. Pyruvate or phosphenalpyruvate formed from glucose by glycolysis may be carboxylated to form malate, which is reduced via fumarate to form succinate. The reduction of fumarate could result in reoxidation of hydrogenated nicotinamide adenine dinucleotide (NADH) under anaerobic conditions, allowing glycolysis and synthesis of adenosine triphosphate (ATP) to be maintained. See CITRIC ACID CYCLE. Donald V. Moore Bibliography. H. G. Bronn (ed.), Klassen und Ordnungen des Tierreichs, vol. 4, 1932–1933; T. Dunagan and D. Miller, Acanthocephala, in Microscopic Anatomy of Invertebrates, vol. 4, pp. 299–332, Wiley, 1991; L. H. Hyman, The Invertebrates: Acanthocephala, Aschelminthes, and Entoprocta, vol. 3, 1951; H. J. Van Cleave, Acanthocephala of North American Mammals, 1953.
Acanthodii A diverse group of early jawed fishes from the Paleozoic Era, usually classified as a distinct class of vertebrates. Most acanthodians have tiny body scales with a cross section like that of an onion, the layers reflecting incremental scale growth. Most acanthodians have strong spines at the leading edges of all fins except for the caudal (tail) fin. Acanthodians mostly were small fishes with large mouths and large eyes, suggesting an active, predatory lifestyle feeding on small prey, including other fishes. The oldest well-preserved examples are of Early Silurian age, and the last surviving examples are of Early Permian age. Primitive forms tended to be from marine environments and to have stouter bodies, larger stouter fin spines, and larger numbers of paired spines in front of the pelvic fins. Later forms are found also in freshwater and estuarine habitats and had more slender bodies with fewer and more slender spines. Although experts do not currently consider acanthodians to be the most primitive of jawed vertebrates (the Placodermi have that distinction), acanthodians are the earliest jawed vertebrates to be represented in the fossil record by complete specimens. Superficially sharklike in appearance, they are instead currently grouped in the Teleostomi as the closest relatives of the bony fishes (Osteichthyes), with which they share characteristics such as small nasal capsules, large eyes, scales capable of continuous growth (in most species), and (in some specialized forms) three pairs of otoliths (ear stones). See OSTEICHTHYES; PLACODERMI; TELEOSTOMI. Anatomy. Most acanthodians are small fishes, with size ranging from a few centimeters to half a meter. A few species reached lengths of more than 2 m (6.6 ft). The body usually was streamlined, with a blunt head, a terminal mouth, and a long, tapered, upturned tail. The large eyes were set close to the front of the head, behind a pair of small nasal capsules. The head usually was covered by large, flattened scales (tesserae), between which ran sensory canals similar to those in bony fishes. The braincase was ossified in advanced forms but not in primitive forms, where it remains cartilaginous. The ear region included three pairs of growing ear stones in some advanced acanthodians; primitive acanthodians had an open connection (endolymphatic duct) between the inner ear and the external environment, through which tiny sand grains entered and were used for the sense of balance, as in many sharks. All acanthodians had jaws (sometimes calcified) consisting of a pair each of upper palatoquadrate cartilages and lower meckelian cartilages. In some acanthodians, jawbones, usually with teeth, were attached to these jaw cartilages (Fig. 1). In many species, teeth were absent; where present, they were of several types: multiple cusps in a whorl attached to a single base, usually near the front of the mouth; separate teeth in whorl-shaped arrangements along the margins of the mouth; or permanent conical teeth and denticles fixed to the jawbones. Gills were positioned behind the head in a compact gill chamber,
39
40
Acanthodii
palatoquadrate cartilage upper jawbone
meckelian cartilage lower jawbone Fig. 1. Jaws of an Early Devonian ischnacanthid from northern Canada. The upper palatoquadrate and lower meckelian cartilages have tooth-bearing jawbones along their facing edges. (From University of Alberta, Laboratory for Vertebrate Paleontology collection)
as in bony fishes; a few acanthodians appear to have had multiple external gill slits on each side, but most had ornamented, platelike bony armor that enclosed the gills, leaving only a single pair of gill slits for water flow. See GILL. Primitive acanthodians had two dorsal fins, an anal fin, and paired pectoral and pelvic fins, all with leading-edge spines (Fig. 2a–c). Later forms retained a single dorsal fin (Fig. 2d). The tail usually was long and flexible. In primitive forms, a series of paired spines, the prepelvic series, was located along the belly in front of the pelvic fins (Fig. 2a). Paired prepectoral spines were present anterior and ventral to the pectoral fins, or were united by bony bases to form a ventral armor between the pectoral fins. These additional paired spines and associated bony plates were reduced or lost in later acanthodians. Scales of acanthodians were small and continuously growing. Most had an onionskin-like structure internally as successive increments of growth were added to the outside of the scale. The crown of the scale was rhombic in outline, made of dentine, and either smooth or ornamented with ridges, while the bulbous base was composed of cellular or acellular bone. Scales of acanthodians often are abundant and useful for dating and correlating rocks. Diversity. At one time there was controversy about whether the most primitive acanthodians were those with many fin spines or those with fewer, but more recently it has been recognized that the most primitive acanthodians were those with multiple pairs of stout, heavily ornamented spines, usually classified in the order Climatiiformes. More advanced acanthodians include those in the orders Ischnacanthiformes and Acanthodiformes. Many other acanthodian species are of uncertain relationships because they are known only from isolated scales or fin spines.
Climatiiformes. This order includes the most primitive acanthodians, typically with broad, heavily ridged spines, multiple pairs of prepelvic spines, and either prepectoral spines or well-developed pectoral armor. Teeth, where present, were numerous and arranged in whorl-like sets superficially similar to those of sharks. Examples are the well-known Climatius, Ptomacanthus, Euthacanthus, Brochoadmones, and Lupopsyrus (Fig. 2a). Diplacanthids, including Diplacanthus, Gladiobranchus, Uraniacanthus, and Tetanopsyrus, sometimes included as a suborder within Climatiiformes, had a pair of large dermal plates on the side of the head, and had their prepelvic spines reduced to one pair or none (Fig. 2b). Climatiiformes appeared in the Silurian; they flourished and then declined in the Devonian. The gyracanthid fishes, sometimes considered to be climatiiform acanthodians, have scales more like those of early chondrichthyans; gyracanthids survived into the Pennsylvanian. Ischnacanthiformes. These acanthodians had tooth whorls at the front of their mouth as well as toothbearing jawbones borne on the upper and lower jaw cartilages (Fig. 1). They had slender fin spines and no prepelvic or prepectoral paired spines. Some ischnacanthiforms were rather large predators. Many species are known only from their toothed jawbones, which often are found separately. Examples include Poracanthodes and Ischnacanthus (Figs. 1, 2c). They are known from the Silurian to the end of the Devonian.
(a)
(b)
(c)
(d) Fig. 2. Examples of acanthodians from the Early Devonian of northern Canada. (a) Primitive climatiiform Lupopsyrus. (b) Diplacanthid Tetanopsyrus. (c) Ischnacanthiform Ischnacanthus. (d) Unnamed mesacanthid acanthodiform. (From the University of Alberta, Laboratory for Vertebrate Paleontology collection)
Acari Acanthodiformes. These are advanced, streamlined, toothless forms with only one dorsal fin and at most one pair of prepelvic spines. One of the best-known acanthodians of any age is Acanthodes, thanks to its ossified internal cranial structures. Later acanthodiforms had long gill rakers and are thought to have been plankton feeders. The group appeared in the Early Devonian and survived until the Early Permian. Examples include Cheiracanthus in Cheiracanthidae, Acanthodes and Homalacanthus in Acanthodidae, and Triazeugacanthus, Melanoacanthus, and Mesacanthus in Mesacanthidae (Fig. 2d). Mark V. H. Wilson Bibliography. R. H. Denison, Acanthodii: Handbook of Paleoichthyology, vol. 5, Gustav Fischer Verlag, Stuttgart, 1979; P. Janvier, Early Vertebrates, Clarendon Press, Oxford, 1996; J. A. Long, The Rise of Fishes: 500 Million Years of Evolution, Johns Hopkins University Press, Baltimore, 1995; J. A. Moy-Thomas and R. S. Miles, Palaeozoic Fishes, 2d ed., W. B. Saunders, Philadelphia, 1971.
Acanthometrida A genus of marine protozoans in the class Acantharea. The kingdom Protozoa contains 18 phyla. One of the parvkingdoms (a hierarchical classification between kingdom and superphylum that is controversial and not officially recognized as such) is the Actinopoda (originally a class) containing two phyla, Heliozoa and Radiozoa. Within the Radiozoa is the class Acantharea, in which the genus Acanthometrida is placed. All members are unicellular planktonic marine protozoa with axopodia. Long slender pseudopodia are found in certain protozoans. Mitochondria are always present and contain flattened cristae. No cilia are present in the trophic phase. Members have a skeleton, made up of strontium sulfate (celestite) limited to 20 radially arranged rods that extend from the center,
forming a characteristic pattern in which angles are quite exact, as in Acanthometra (see illustration). Since strontium sulfate is soluble in seawater, no acantharean fossils exist. The central endoplasm contains numerous nuclei and other organelles. The cytoplasm surrounding the spines contains a conical array of contractile microfilaments (myophrisks or myonemes that are Ca2+-activated) containing the contractile protein actin. The microfilaments expand or contract the gelatinous sheath surrounding the cell, a phenomenon apparently responsible for changes in level of flotation. Electron microscopic research with related genera shows that the microfilaments are arranged in two sets of bands: one set runs parallel to the axis of the spine, and the other runs crosswise, thus producing a crossfibrillar network. Cysts and flagellated swarmers represent two stages of the life cycle, which are not completely known at present. The endoplasm often contains zooxanthellae (small dinoflagellates living in the protozoan’s endoplasm). See ACANTHAREA; ACTINOPODEA; CELESTITE; PROTOZOA; SARCODINA; SARCOMASTIGOPHORA; STRONTIUM. O. Roger Anderson; John P. Harley Bibliography. T. Cavalier-Smith, Kingdom Protozoa and its 18 phyla, Microbiol. Rev., 57(4):953– 994, 1993; M. Grell, in H.-E. Gruner (ed.), Lehrbuch der Speziellen Zoologie, 4th ed., Gustav Fischer, Stuttgart, 1980; K. Hausmann and N. Hulsmann, Protozoology, Georg Thieme Verlag, New York, 1996.
Acanthophractida An order of Acantharea. In this group of protozoa, skeletons typically include a latticework shell, although the characteristic skeletal rods are recognizable. The latticework may be spherical or ovoid, is fused with the skeletal rods, and is typically concentric with the central capsule. The body is usually covered with a single or double gelatinous sheath through which the skeletal rods emerge. Myonemes (usually a specific number) extend from the gelatinous sheath to each skeletal rod. These marine forms live mostly below depths of 150– 200 ft (45–60 m). The order includes Coleaspis, Diploconus, Dorotaspis, and many other genera. See ACANTHAREA; ACTINOPODEA; PROTOZOA; SARCODINA; SARCOMASTIGOPHORA. Richard P. Hall
central capsule central capsule membrane extracapsular protoplasm zooxanthellae
skeletal spicule myophrisks
Acanthometra. (After L. H. Hyman, The Invertebrates, vol. 1, McGraw-Hill, 1940)
Acari A subclass of Arachnida, the mites and ticks. All are small (0.004–1.2 in. or 0.1–30 mm, most less than 0.08 in. or 2 mm in length), have lost most traces of external body segmentation, and have the mouthparts borne on a discrete body region, the gnathosoma (Fig. 1). They are apparently most closely related to Opiliones and Ricinulei. General characteristics. The chelicerae may be chelate or needlelike. Pedipalps are generally smaller than the walking legs, and are simple or sometimes
41
42
Acari
gnathosoma
Fig. 1. Ventral view of Laelaps echidnina, the spiny rat mite (Parasitiformes), showing major body divisions.
weakly chelate. Legs typically move in the slow walking gait of arachnids; they may, however, be modified by projections, enlarged claws, heavy spines, or bladelike rows of long setae for crawling, clinging to hosts, transferring spermatophores, or swimming. Mites have a variety of integumentary mechanoreceptors and chemoreceptors (often modified setae), and one or two simple eyes (ocelli) frequently occur anterolaterally on the idiosoma (occasionally one anteromedially as well). The ganglia of the central nervous system are coalesced into a single “brain” lying around the esophagus. Excretion is by guanine crystals secreted into the hindgut and voided through a ventral anus. Respiration is by a tracheal network fed by one to four pairs of spiracles that vary in position. A simple dorsal heart is present in a few larger forms. Life cycle. The typical acarine life cycle consists of six instars: egg; six-legged larva (Fig. 2); three nymphal instars (protonymph, deutonymph, and tritonymph); and adult, usually similar to active
idiosoma
gnathosoma
Fig. 2. Ventral view of a Hygrobateslarva (Acariformes) modified for aquatic life.
nymphs. Variations include suppression of the larva, one or more inactive nymphs for internal structural reorganization or phoresy, and change in number of nymphal stages. Male spermatophores may be deposited on the substrate, transferred to the female gonopore by modified legs or chelicerae, or inserted directly by a male copulatory organ. Females are normally oviparous, depositing 4–10,000 eggs singly or in clusters. The generation time varies from 4 days for some terrestrial forms to as much as 2 years in some fresh-water species. Ecology and feeding. Acarine habitats and feeding habits are much more diverse than in other arachnids. Mites are ubiquitous, occurring from oceanic trenches below 13,000 ft (4000 m) to over 21,000 ft (6300 m) in the Himalayas and suspended above 3300 ft (1000 m) in the atmosphere; there are species in Antarctica. Soil mites are taxonomically varied, show the least specialized adaptations to habitat, and are frequently well-sclerotized predators. Inhabitants of stored grains, cheese, and house dust (many causing allergic reactions in humans) are also relatively unspecialized. The dominant forms inhabiting mosses are heavily sclerotized beetle mites (Oribatei), some of which can curl into a ball to reduce moisture loss from the ventral surface. Flowering plant associates include spider mites (Tetranychidae), which suck the contents of leaf epidermal cells, and gall mites (Eriophyidae), which cause a neoplastic growth of host tissue that provides both shelter and food. Fungivores are usually weakly sclerotized inhabitants of moist or semiaquatic habitats, while the characteristic fresh-water mites include sluggish crawlers, rapid swimmers, and planktonic drifters, all with reduced idiosomal setation. Mites associated with invertebrate animals include internal, external, and social parasites as well as inactive phoretic stages of Insecta, Crustacea, Myriapoda, Chelicerata, Mollusca, and Parazoa. A similar diversity of species are parasites or commensals of vertebrates; microhabitats include the nasal mucosa and lungs of reptiles and mammals, burrows in the skin of humans and most other vertebrates (scabies and mange mites), hair follicles (Demodicidae), feathers (both externally and within the feather shaft), fur (bat mites), and host nests and burrows. Some parasitic mites are disease vectors. Origin and classification. Mites appear in the Devonian among the first true arachnids and thus apparently diverged early from the basic chelicerate stock. Differences in optical activity of the setae and degree of idiosomal segmentation suggest the possibility of a diphyletic origin. Over 30,000 species have been described, and it is estimated that as many as 500,000 may exist. See ARACHNIDA. David Barr Bibliography. D. R. Cook, Water Mite Genera and Subgenera, Mem. Amer. Entomol. Inst. 21, 1974; A. Kaestner, Invertebrate Zoology, vol. 2: Arthropod Relatives, Chelicerata, Myriapoda, 1968; G. W. Krantz, A Manual of Acarology, 2d ed., 1978; S. P. Parker (ed.), Synopsis and Classification of Living Organisms, 2 vols., 1982; A. E. Treat, Mites of Moths and Butterflies, 1975.
Accelerating universe
Accelerating universe Cosmic expansion is speeding up. In 1998, astronomers showed evidence from exploding stars that the expansion of the universe is not slowing down because of the gravity of matter in the universe. By observing the light from distant supernovae, astronomers can measure the distance to each explosion. By studying the redshift in the supernova’s spectrum, they can measure the expansion of the universe since the light was emitted and, from a set of these measurements, can trace the history of cosmic expansion. Instead of seeing a slowing universe, as expected, observations indicate there is a component of the universe, called the dark energy, that makes cosmic expansion speed up. Dark energy comprises about two-thirds of the energy density of the universe. While it may be related to Albert Einstein’s controversial cosmological constant, its true nature remains to be discovered by observation and thought. See COSMOLOGICAL CONSTANT; DARK ENERGY; SUPERNOVA. Astronomers have known since the 1920s that galaxies are separated by millions of light-years and that the space between them is stretching, moving the galaxies apart: the nearest galaxies are moving away from us slowly and distant galaxies are receding more rapidly. This relation was discovered empirically by Edwin Hubble using the 100-in. (2.5-m) telescope at the Mount Wilson Observatory. If our patch of the universe is not unique, the simplest explanation for this observation is that the universe as a whole is expanding in all directions. What is new about the observations made in the 1990s is that the distances and times probed by the supernovae are large enough for the effects of cosmic acceleration to be measurable. Since light travels at a finite speed, the history of cosmic expansion can be traced by looking at distant objects, whose light left their source hundreds of millions or even billions of years in the past. This procedure makes a telescope into a time machine capable of showing how the expansion of the universe has changed over time. Einstein’s general theory of relativity, applied in the cosmic setting, shows that the presence of matter in the universe should lead to the gradual slowing of cosmic expansion. Since the 1950s, astronomers had sought to measure this deceleration of the universe, hoping to use it as a guide to the density of matter in the universe and as a predictor of the future behavior of cosmic expansion—whether the universe would coast outward indefinitely or would be drawn back to a dense state like the big bang from which it emerged 14 × 109 years ago. See BIG BANG THEORY; RELATIVITY. In the 1990s, the perfection of accurate distancemeasuring techniques based on supernova explosions and the development of powerful observing techniques made it possible to carry through the astronomers’ quest to measure changes in the cosmic expansion rate. Instead of finding that the expansion has slowed over time because of the decelerating effects of gravity, as most expected, the new results
from two teams showed the opposite: we live in an accelerating universe. Contrary to all expectation, the universe’s expansion is speeding up. Subsequent work on supernovae, on the clustering of galaxies, and on the faint cosmic background light left over from the big bang itself converge on the same solution: the universe today is dominated by a strange dark energy whose role in fundamental physics is certainly important, but of which our present understanding is very limited. See COSMIC BACKGROUND RADIATION. Expanding universe. Before 1920, astronomers thought that the Milky Way Galaxy, of which the Sun is an inconspicuous member, constituted the entire universe. When Einstein applied his recently developed theory of gravitation, the general theory of relativity, to the universe in 1916, he was guided by Willem DeSitter, a leading astronomer, to think of a static system. Einstein found that a static solution to his equations could be constructed if he added in a term that we now call the cosmological constant. This term had the effect of balancing out the attractive force of gravity with a kind of cosmic repulsion to produce a static universe. Hubble, working at the Mount Wilson Observatory, showed that the distances to what he called the spiral nebulae were much larger than the extent of the Milky Way, and that it was more correct to think of these as independent systems, galaxies, each equivalent to the Milky Way. He did this by measuring the apparent brightness of stars whose intrinsic brightness he knew—Cepheid variables, whose intrinsic brightness can be found by measuring their periods of variation. The fainter a star appears, for a given intrinsic brightness, the more distant it must be. By observing this, a particular type of star, Hubble was able to show that the distances to even the nearest nebulae were a few million light-years, far larger than the size of the Milky Way, whose own size also had to be correctly measured. See CEPHEID; GALAXY, EXTERNAL; MILKY WAY GALAXY. Hubble was also able to use measurements of the spectra of galaxies to show that their spectra have a redshift, the signature that they are receding from us. By combining the measurements of distance and of recession, Hubble discovered that the more distant galaxies are moving away from us more rapidly. The modern interpretation of these observations is cosmic expansion, with the universe stretching out in all directions at a rate of about 1 part in 14,000,000,000 of its size each year. See HUBBLE CONSTANT; REDSHIFT. Naturally, Einstein had to adjust to this new picture of the cosmos. He had inserted his cosmological constant to make a static world, but later observations showed that the world was not static. Einstein was then inclined to banish the cosmological constant, which he knew was not mathematically justified and which he never did like much anyway. But DeSitter was not so sure this was a good idea. He thought it was possible that the universe had started out slowly but had accelerated as a result of the repulsive properties of the cosmological constant, denoted by
43
Accelerating universe
Fig. 1. Supernova, 1994D. This type Ia supernova is in a galaxy at a distance of about 5 × 107 light-years in the Virgo cluster of galaxies. For a month, the light from a single exploding white dwarf is as bright as 4 × 109 stars like the Sun. (P. Challis, Center for Astrophysics/STScI/NASA)
(lambda). By 1932, Einstein and DeSitter were ready to stop discussing , leaving that for future generations. In a joint paper they said of the cosmological constant, “An increase in the precision of data. . . will enable us in the future to fix its sign and determine its value.” Type Ia supernovae. That future arrived in 1998, 66 years later. The measurements on which the evidence for cosmic acceleration rests are similar to those carried out by Hubble: except that the stars being used are a million times brighter and the distances to the galaxies in which they are located are a thousand times larger. The claim that the universe is accelerating is so extraordinary that it is worth examining the evidence on which it is based: the brightness of exploding stars astronomers call type Ia supernovae (Fig. 1). Type Ia supernovae are a particular type of exploding star that can be recognized by its spectrum. They are found in all types of galaxies and are thought to be the result of a sudden thermonuclear burning wave that propagates through a white dwarf star. This flame burns the carbon and oxygen of the star’s interior into radioactive nickel. The subsequent decay into cobalt and iron releases energy gradually as the star is blown apart. For about a month, a type Ia supernova shines as brightly as 4 × 109 stars like the Sun shining simultaneously. This property makes them visible over very large distances, large enough to detect the change in cosmic expansion rate that has taken place while the light was on its way from a type Ia supernova halfway across the universe. But being bright is only half the story. In order to measure cosmic distances from apparent brightness with good precision, the ideal tool would be a “standard candle” for which all the objects had the same intrinsic brightness. Then a measurement
of the apparent brightness would give the distance. Type Ia supernovae are pretty good standard candles: the intrinsic light output at their peak luminosity varies by only a factor of 3 or so. This uniform behavior probably is connected to the theoretical prediction that there is a well-defined upper mass limit for white dwarfs of about 1.4 solar masses that is set by the quantum-mechanical properties of electrons. See WHITE DWARF STAR. Even better precision comes from empirical approaches to type Ia supernovae. It turns out that the brightest supernovae have distinctly slower declines in brightness after they reach their peak, while the dimmer type Ia supernovae decline more rapidly. By measuring the rise and fall in the brightness of a supernova, which does not depend on its distance, one can tell whether an extrabright supernova or an extradim one is being observed, and take that information into account in estimating the distance. When the shape of the light curve is taken into account, the uncertainty in the distance to a single supernova shrinks to under 10%, making this type of supernova the very best distance-measuring tool for surveying the universe (Fig. 2). See LIGHT CURVES. Technological advances. The disadvantage of type Ia supernovae is that they are rare. In a typical large galaxy, the rate of type Ia supernova explosions is about one in a century. In our own Milky Way Galaxy, the last event that was probably a type Ia supernova was observed by Tycho Brahe in 1572. So if you want to find and measure many supernovae to survey the history of cosmic expansion, you need either to be very patient or to be able to search many galaxies for the desired events. This point is where technology has made great strides. The electronic light detectors used in modern telescopes are charge-coupled devices (CCDs), similar to those found in digital cameras. These are 4
velocity, 104 km s−1
44
3
2
1
0
0
100
200 300 400 distance, Mpc
500
600
Fig. 2. Modern Hubble diagram for 95 type Ia supernovae. Recession velocity is plotted against distance in megaparsecs (1 Mpc = 3.26 × 106 light-years = 3.09 × 1019 km). The points represent supernovae up to 2 × 109 light-years away, with recession velocities up to one-tenth of the speed of light. The small scatter about the diagonal line shows that type Ia supernovae make excellent standard candles for measuring distances in the universe. This nearby patch does not show the effects of acceleration. (Based on data from Saurabh Jha)
Accelerating universe more than 100 times more sensitive to light than the photographic plates used in Hubble’s time. As the technology for fabricating computer chips has improved, the CCDs available to astronomers have grown larger as well as more sensitive, so that now arrays of 1000 megapixels are being applied to supernova studies and still larger cameras are in devlopment. These large and sensitive devices allow a modern telescope to image many thousands of galaxies in a single exposure. By coming back night after night to the same part of the sky, supernova searches in a single night gather the data needed to find many of these rare events (Fig. 3). See ASTRONOMICAL IMAGING; CHARGE-COUPLED DEVICE. But gathering the light is only part of the story of finding supernovae to do cosmology. You also need to detect the new starlike events among the many fuzzy galaxies that do not change from night to night. This is the place where advances in computer hardware and software have turned the task of supernova searching from handicraft into something on an industrial scale. By using a fast computer to subtract digital images obtained last month (or last year) from the ones obtained tonight, it is relatively straightforward to sift out overnight a handful of new objects from the hundreds of thousands that remain unchanged. Measuring the apparent brightness of these supernovae and the shape of their rise and fall over a month gives enough information to measure the distance to a supernova. Measuring the spectrum of the supernovae can show whether it is of type Ia and measure the redshift. By measuring the redshift at various distances, one can trace the history of cosmic expansion to look for the effects of slowing that results from gravity or acceleration caused by dark energy. The first attempt at this type of searching was in 1988 by a Danish group. Because only small CCD chips were available, they published data on only one type Ia supernova. But the Supernova Cosmology Project based at Lawrence Berkeley Laboratory developed effective search techniques, and by 1997 they had a preliminary result. Their first indication, based on seven supernovae found in 1994 and 1995, was that the universe was slowing down because of the braking effects of dark matter. In 1998, the High-Z Supernova Search Team published a result indicating the opposite: that the universe was accelerating. Results from the Supernova Cosmology Project published in 1999 showed the same thing: we live in a universe where the rate of cosmic expansion is increasing as the universe grows older. Subsequent observations have confirmed this latter result, and both teams agree. The samples have grown larger, the span of cosmic time probed by the supernova observations has increased, and methods for dealing with the pernicious effects of absorption by dust have been developed. One especially interesting result comes from a supernova search and follow-up carried out using the Hubble Space Telescope both to find and to get the light curves and spectra of very distant supernovae. This work
epoch 1
epoch 2
45
epoch 2 − epoch 1
Fig. 3. Subtracting images to find supernovae. The image from a month ago (epoch 1) is subtracted from last night’s image (epoch 2) to reveal a new supernova. The image area shown is about 1/1000 of the full area provided by the CCD camera. Rapid processing of dozens of image pairs demands nimble software and a dedicated team of searchers. (Brian Schmidt, Australian National University)
showed that cosmic acceleration is a comparatively recent phenomenon, acting over only the last 5 × 109 years or so. Supernovae discovered at larger distances than 5 × 109 light-years show that the universe was probably slowing down when it was younger. So there is evidence that the universe was slowing down when it was young and dense, presumably resulting from the presence of dark matter, but then there came a shift, about 9 × 109 years after the big bang (and 5 × 109 years before the present), from deceleration to acceleration. In the lexicon of elementary physics, change in position is velocity, change in velocity is acceleration, and change in acceleration is called jerk. This program with the Hubble Space Telescope has produced evidence for the cosmic jerk. See DARK MATTER; HUBBLE SPACE TELESCOPE. Dark energy?. By combining measurements of supernova distances with results from galaxy clustering and from fluctuations in the cosmic microwave background, consensus has emerged that the universe at present has about 73% of its energy density in dark energy, about 23% in dark matter, and only 4% in the form of atoms made of neutrons, protons, and electrons. While it is satisfying that there is a single solution that matches a wide variety of current data, an honest appraisal is that we do not know what the dark energy is or what the dark matter is. The dark energy could be a form of Einstein’s cosmological constant. The cosmological constant has the right property: a tendency to make the universe expand faster over time. This tendency is why DeSitter thought the cosmological constant might be responsible for the cosmic expansion observed by Hubble. But the modern view of the cosmological constant is that it has its origin in the quantum properties of the vacuum. When the same principles that work so well for electromagnetism are applied to gravitation, the result is a quantitative disaster of unprecedented proportions. The observed dark energy is about 10120 times smaller than the predicted amount.
46
Acceleration Another possibility is that the dark energy is not something that is constant in time, as the cosmological constant would be, but something that has been changing its value as the universe unfolds. This property could be detected by observations of distant supernovae. The next round of observational programs for distant supernovae is aimed at observing the history of cosmic expansion with enough precision to tell the difference between a dark energy that does not change (like the cosmological constant) or something that varies with time. Large samples of supernovae at moderate distances and moderate samples of supernovae at large distances are being constructed. The goal is to pin down the properties of the dark energy well enough to tell if it is different in any way from a cosmological constant with a very small numerical value. So far, no significant deviation has been found, but the samples published to date are small. Dark energy is so new, so important, and so mysterious that it will be the focus of strenuous efforts in observational astronomy for decades to come. See COSMOLOGY; UNIVERSE. Robert P. Kirshner Bibliography. D. Goldsmith, The Runaway Universe, Perseus Books, Cambridge, MA, 2000; R. P. Kirshner, The Extravagant Universe: Exploding Stars, Dark Energy, and the Accelerating Cosmos, Princeton University Press, 2002; L. Krauss, Quintessence, Basic Books, New York, 2000.
a constant direction, the acting force will produce a continuous change in speed. If it is moving with a constant speed, the acting force will produce an acceleration consisting of a continuous change of direction. In the general case, the acting force may produce both a change of speed and a change of direction. R. D. Rusk Angular acceleration. This is a vector quantity representing the rate of change of angular velocity of a body experiencing rotational motion. If, for example, at an instant t1, a rigid body is rotating about an axis with an angular velocity ω 1, and at a later time t2, it has an angular velocity ω 2, the average angular acceleration α is given by Eq. (2), expressed α=
ω2 − ω1 ω = t2 − t1 t
(2)
in radians per second per second. The instantaneous angular acceleration is given by α = dω/dt. Consequently, if a rigid body is rotating about a fixed axis with an angular acceleration of magnitude α and an angular speed of ω0 at a given time, then at a later time t the angular speed is given by Eq. (3). ω = ω0 + αt
(3)
A simple calculation shows that the angular distance θ traversed in this time is expressed by Eq. (4). ω0 + (ω0 + αt) θ = ωt = t 2 (4) = ω0 t + 1 /2 αt 2
Acceleration The time rate of change of velocity. Since velocity is a directed or vector quantity involving both magnitude and direction, a velocity may change by a change of magnitude (speed) or by a change of direction or both. It follows that acceleration is also a directed, or vector, quantity. If the magnitude of the velocity of a body changes from υ 1 ft/s to υ 2 ft/s in t seconds, then the average acceleration a has a magnitude given by Eq. (1). To designate it fully the direction should be a=
velocity change v2 − v1 v = = elapsed time t2 − t1 t
(1)
given, as well as the magnitude. See VELOCITY. Instantaneous acceleration is defined as the limit of the ratio of the velocity change to the elapsed time as the time interval approaches zero. When the acceleration is constant, the average acceleration and the instantaneous acceleration are equal. If a body, moving along a straight line, is accelerated from a speed of 10 to 90 ft/s (3 to 27 m/s) in 4 s, then the average change in speed per second is (90 − 10)/4 = 20 ft/s or (27 − 3)/4 = 6 m/s in each second. This is written 20 ft per second per second or 20 ft/s2, or 6 m per second per second or 6 m/s2. Accelerations are commonly expressed in feet per second per second, meters per second per second, or in any similar units. Whenever a body is acted upon by an unbalanced force, it will undergo acceleration. If it is moving in
In the illustration a particle is shown moving in a circular path of radius R about a fixed axis through O with an angular velocity of ω radians/s and an angular acceleration of α radians/s2. This particle is subject to a linear acceleration which, at any instant, may be considered to be composed of two components: a radial component ar and a tangential component at. Radial acceleration. When a body moves in a circular path with constant linear speed at each point in its path, it is also being constantly accelerated toward the center of the circle under the action of the force required to constrain it to move in its circular path. This acceleration toward the center of path is called radial acceleration. In the illustration the radial acceleration, sometimes called centripetal acceleration, is shown by the vector ar. The magnitude of its value is υ 2/R, or ω2/R, where υ is the instantaneous linear
at
ar
O
Radial and tangential accelerations in circular motion.
Acceleration analysis velocity. This centrally directed acceleration is necessary to keep the particle moving in a circular path. Tangential acceleration. The component of linear acceleration tangent to the path of a particle subject to an angular acceleration about the axis of rotation is called tangential acceleration. In the illustration, the tangential acceleration is shown by the vector at. The magnitude of its value is αR. See ACCELERATION MEASUREMENT; ROTATIONAL MOTION. C. E. Howe; R. J. Stephenson
1 4
P
D 1
ω2
O
3
t
V PB
VB
2
VB
⊥ DP (a)
(b)
B
Acceleration analysis A mathematical technique, often done graphically, by which accelerations of parts of a mechanism are determined. In high-speed machines, particularly those that include cam mechanisms, inertial forces may be the controlling factor in the design of members. An acceleration analysis, based upon velocity analysis, must therefore precede a force analysis. Maximum accelerations are of particular interest to the designer. Although an analytical solution might be preferable if exact maxima were required, graphical solutions formerly tended to be simpler and to facilitate visualization of relationships. Today, for advanced problems certainly, a computer solution, possibly one based on graphical techniques, is often more effective and can also produce very accurate graphical output. See FORCE; KINEMATICS; MECHANISM. Accelerations on a rigid link. On link OB (Fig. 1a) acceleration of point B with respect to point O is the vector sum of two accelerations: (1) normal acceleration AnBO of B with respect to O because of displacement of B along a path whose instantaneous t
A BO
B α t
V BO ω
O (a)
B
OB
A BO n
⊥ OB
ABO
(b) Fig. 1. Elementary condition. (a) Point B on rigid member rotates about center O. (b) Vector diagram shows normal, tangential, and resultant accelerations.
A nB PB A nPB
2
V B OB OB
2
V PB PB ⊥ PB
AP
⊥ PD
t
direction of A PB
(c)
⊥ PB
t
VP
PD t
direction of A P
A
n P
OA V
2 P
PD
Fig. 2. Four-bar linkage. (a) Given are the linkages 2, 3, and 4 between fixed supports 1–1 and the velocity of point B, ω 2 = 0. (b) A vector polygon is used to determine the velocity at point P. (c) A graphic solution then finds acceleration at P.
center of curvature is at O, and (2) tangential acceleration of AtBO of B with respect to O because of angular acceleration α. For conditions of Fig. 1a with link OB rotating about O at angular velocity ω and angular acceleration α, the accelerations can be written as Eqs. (1) and (2). The vector sum or resultant is ABO (Fig. 1b). AnBO = (OB)ω2 = VB2 /(OB)
(1)
AtBO = (OB)α
(2)
Accelerations in a linkage. Consider the acceleration of point P on a four-bar linkage (Fig. 2a) with α 2 = 0 and hence ω2 = k, the angular velocity of input link 2. First, the velocity problem is solved yielding VB and, by Fig. 2b, VPB. Two equations can be written for AP; they are solved simultaneously by graphical means in Fig. 2c by using Fig. 2b; that is, normal accelerations of P with respect to B and D are computed first. Directions of tangential acceleration vectors AnPB and AnP are also known from the initial geometry. The tip of vector AP must lie at their intersection, as shown by the construction of Fig. 2c. See FOUR-BAR LINKAGE. Explicitly, acceleration AP of point P is found on the one hand by beginning with acceleration AB of point B, represented by Eq. (3). To this acceleration Ap = AB → APB
(3)
47
48
Acceleration analysis n is added vectorially normal acceleration APB and tant gential acceleration APB, which can be written as Eq. (4). Also for link 2, α 2 = 0 and ABt = 0; and AB need not be split up.
VP4P2
VP2
1
VP4
(4)
AP = AB →
AnPB
→
AtPB
follower
(4) P
2
On the other hand, AP can also be expressed as Eq. (5), or in Fig. 2c. The intersection of the two tan-
(1) cam
AP =
AnP
→
Atp
ω2
O
(5)
(3)
gential components defines the tip of AP. This problem illustrates the generalization that any basic acceleration analysis can be thoroughly performed only after completion of the underlying velocity analysis. Acceleration field. Acceleration is a vector and hence has magnitude and direction but not a unique position. A plane rigid body, such as the cam of Fig. 3a, in motion parallel to its plane will possess an acceleration field in that plane differing from, but resembling, a rigid-body velocity vector field. See CAM MECHANISM. Every acceleration vector in this field, absolute or relative, makes with its instantaneous radius vector the same angle γ (tangent γ = α/ω2) and is proportional in magnitude to the length of this radius vector (Fig. 3b). The acceleration field at any instant is thus defined by four parameters: magnitude and direction of accelerations at any two points on the body. From an acceleration field defined by aA and aB (Fig. 3c), one can quickly determine the instantaneous center of zero acceleration. From point A, construct vector AA parallel and equal to acceleration aB, represented by vector BB. The relative acceleration aAB is given by the vector of Eq. (6).
(2)
(a)
(B)
r (α2 + ω4 )1/2
r γ
C
AA Ac O
(b)
B
γ
A'
γ
AB
& = B'B
Γ
B'
aA
A'' aAB
γ
aB B
(6)
Angle γ between resultant aAB and line AB is common to all radii vectors; therefore, construct radii vectors through A and through B both at γ (properly signed) to their respective acceleration vectors. These lines intersect at , the center of zero acceleration. The geometry contains two similar triangles; AB is similar to AAA because corresponding sides are equally inclined by the angle γ . From these similar triangles comes the generalization that acceleration magnitudes are proportional to radii vectors, just as for the velocity field. Angle γ (tan γ = α/ω2) is also used in orientation of acceleration image polygons, especially if the preceding velocity analysis used the image polygon method. Points not on the same link. Acceleration of follower link 4 of the cam mechanism (Fig. 3a) is determined by considering relative accelerations of coincident points not on the same link. Thus point P on link 4, designated P4, and coincident point P on link 2, designated P2, have instantaneous velocities, as shown by vectors VP4 and VP2 (Fig. 3a). Relative velocity of P4 with respect to P2 is VP4P2 as indicated. Accelera-
rω4
γ
A
aAB = aA − aB
rα
A
γ (c) Fig. 3. Accelerations on a cam and follower. (a) The cam mechanism. (b) The acceleration field about a rigid rotating structure. (c) An acceleration field determining an instantaneous center of rotation.
tion of P4 can now be determined by Eq. (7), where AP4 = AP2 → AnP4 P2 → AtP4 P2 → 2ω2 VP4 P2
(7)
the last term is the Coriolis component. This component results from referring motion of P4 on body 4 to P2 on body 2. Serious errors of analysis can result from omission of the Coriolis component. See CORIOLIS ACCELERATION. Occasionally, mechanisms that appear to require the analysis described in the preceding paragraph may also be analyzed by constructing an instantaneously equivalent linkage shown by broken lines (2), (3), and (4) in Fig. 3a. The instantaneous equivalent linkage is then analyzed by the method of Fig. 2. Alternative procedures. Other methods may also be used. It is not always necessary to resolve vectors
Acceleration measurement into normal and tangential components; they can be solved by conventional orthogonal component techniques. For points A and B in part m, Eq. (8) holds. AB AB AB 2 aAB A = aB + aAB = aB + ω · AB
(8)
A convenient approximate solution is obtained by considering the difference between velocities due to a small angular displacement. This difference in velocities divided by the elapsed short time interval approximates the acceleration. Typically, the angular displacement is taken to be 0.1 radian (0.05 radian on either side of the position under study). See ACCELERATION; VECTOR (MATHEMATICS). Douglas P. Adams Bibliography. J. Angeles, Spatial Kinematics Chains: Analysis-Synthesis-Optimization, 1982; J. S. Beggs, Kinematics, 1983; C. W. Ham et al., Mechanics of Machinery, 4th ed., 1958; D. Lent, Analysis and Design of Mechanisms, 2d ed., 1970; A. Sloane, Engineering Kinematics, 1966.
Acceleration measurement The technique of measuring the magnitude and direction of acceleration. Measurement of acceleration includes translational acceleration with tangential and radial components and angular acceleration. See ACCELERATION. Translation acceleration. The components of translational acceleration are the tangential one due to a change in magnitude and the normal one due to a change in direction of the velocity vector. Acceleration measurement by instruments is commonly expressed in terms of g, a dimensionless ratio of acceleration divided by the local value of gravity. When all points of a rigid body are traveling along parallel straight line paths, the body is said to be in linear motion. For such motion the normal component of acceleration is absent, and if the velocity of the body is changing, the resulting translational acceleration is called linear acceleration. There are many familiar applications of linear acceleration measurement, for example, in land vehicles such as automobiles, in aerospace vehicles, and in moving machinery components. One of the purposes of acceleration measurement is to determine the force causing the acceleration as manifested by Newton’s second law. Also, from acceleration versus time data it is possible by either graphical or numerical integrations to determine the corresponding velocity-time and then the displacement-time histories. Since such techniques are comparatively accurate, it is evident that acceleration measurement is also important in the determination of velocity and displacement. An acceleration-measuring instrument, or accelerometer, essentially consists of a reference mass mounted by means of a suitable configuration within a housing. A triaxial accelerometer will have three such reference masses so that they respond individually in three mutually perpendicular directions. Ac-
celerometers with a single reference mass respond to accelerations along a predefined axis. To measure translational acceleration whose exact direction is unknown or whose direction varies, either a single triaxial accelerometer or three ordinary accelerometers are required. See ACCELEROMETER. An indirect determination of acceleration is possible by measuring the displacement-time history of a system undergoing linear acceleration. A graphical or a numerical determination of the slope or derivative at selected values of time will produce the corresponding velocity-time relationship, and a repetition of the procedure yields the acceleration history. The accuracy of this means of obtaining acceleration values is poor because the second set of slope determinations by such approximate methods is inherently inaccurate. The acceleration g due to gravity on the Earth can be determined by various means. A reasonably accurate means of determining g at a given location is by accurately measuring the period T and length L of a simple pendulum on a long cord and subject to a small oscillation in still air. The value of g isthen readily determined from the formula T = 2π L/g. Another means of determining g is offered by measuring the deflection of a reference mass supported by a highly accurate spring with known characteristics. Angular acceleration. Technically, a body is said to have angular motion if some fixed reference line between two points rotates or sweeps out an angle. With this definition it is evident that a body can travel along a curved path and still not rotate. For example, the chairs on a moving ferris wheel are not in angular motion with respect to ground. See ANGULAR MOMENTUM. The measurement of angular acceleration is useful in applications such as the angular motion of aerospace vehicles and the components of rotating machinery. Since angular and translational motion are kinematically related, it is possible, when the exact location of the axis of rotation is known, to determine angular acceleration by measuring the corresponding tangential component of translational acceleration. For determining angular acceleration, a suitably suspended reference mass responds to an angular acceleration about an axis parallel to, or coincident with, its own predefined axis of rotation. As is the case for the translational acceleration measuring instruments, various designs of the angular versions are possible. A direct means of measuring angular acceleration is afforded by the angular accelerometer. An indirect means of determining angular acceleration is to measure the angular velocity of a body. This technique is applicable in situations where a body is in pure rotation about a fixed axis. The direction of the rotating body can be determined from the output of an electric generator connected to it on the same shaft. A stroboscope can also be used for the same purpose. The rotational speed versus time history can then be graphically or numerically
49
50
Accelerator mass spectrometry differentiated once to yield the angular acceleration. See ROTATIONAL MOTION; STROBOSCOPE. Teruo Ishihara; Roger C. Duffield Bibliography. T. G. Beckwith, R. D. Marangoni, and J. H. Lienhard, Mechanical Measurements, 6th ed., 2006; R. S. Figliola and D. E. Beasley, Theory and Design for Mechanical Measurements, 4th ed., 2006; R. C. Hibbeler, Engineering Mechanics: Statics and Dynamics, 10th ed., 2003; W. F. Riley and L. D. Sturges, Engineering Mechanics: Dynamics, 2d ed., 1995.
Accelerator mass spectrometry The use of a combination of mass spectrometers and an accelerator to measure the natural abundances of very rare radioactive isotopes. These abundances are frequently lower than parts per trillion. The most important applications of accelerator mass spectrometry are in archeological, geophysical, environmental, and biological studies, such as in radiocarbon dating by the counting of the rare carbon-14 (radiocarbon; 14 C) isotope. See MASS SPECTROSCOPE; PARTICLE ACCELERATOR. The advantage of counting the radioactive atoms themselves rather than their decay products is well illustrated by radiocarbon dating, which requires the measurement of the number of 14C atoms in a sample. This number can be inferred from the beta-particle
Low-energy mass spectrometry
emission rate or by counting directly the number of radioactive atoms. The long half-life of 5730 years for 14 C implies that only 15 beta-particle emissions per minute are observed from 1 g (0.035 oz) of contemporary carbon which contains 6 × 1010 atoms of 14C. Conventional mass spectrometry has not been able to accomplish the task of counting the 14C atoms directly. However, an accelerator mass spectrometer can be used to count the 14C atoms at over 30 per second from a milligram sample of carbon. Consequently, accelerator mass spectrometry can be used to date samples that are a thousand times smaller than those that are dated by using the beta-particle counting method, and the procedure is carried out about 120 times faster. The advantage of accelerator mass spectrometry over beta-particle counting is even greater for very long-lived radioactive isotopes such as aluminum-26 (7 × 105 year half-life) and iodine-129 (1.6 × 107 year half-life). See RADIOACTIVITY; RADIOCARBON DATING. Apparatus. The success of accelerator mass spectrometry results from the use of more than one stage of mass spectrometry and at least two stages of ion acceleration. The layout of an ideal accelerator mass spectrometer for radiocarbon studies is generally divided for convenience into three stages (see illus.). In the first, low-energy, mass spectrometry stage, there is the usual acceleration at the negative-ion source to produce a directed beam of ions with
High-energy mass spectrometry
Tandem accelerator insulating gas
voltmeter
A2
magnet
A3 L3
L4
F1 magnet box A1
magnet
electrostatic analyzer F2
acceleration tube
106 -volt electrode
stripping canal
C-12 C-13
A4
ion currents
electrostatic analyzer
electrostatic analyzer
L2 sample holder cesium gun C− ion source
L1
C-14 ion detector
Simplified diagram of an accelerator mass spectrometer used for radiocarbon dating. The equipment is divided into three sections. Electric lenses L1–L4 are used to focus the ion beams. Apertures A1–A4 and charge collection cups F1 and F2 are used for setting up the equipment. The cesium ions from the gun create the negative ions of carbon at the surface of the sample.
Accelerator mass spectrometry an energy of several kiloelectronvolts. This result is achieved by raising the sample holder to the required negative voltage. The first acceleration is followed by an electrostatic lens and an electrostatic analyzer for further ion focusing and energy selection, so that the subsequent mass analysis by a magnetic analyzer can select, uniquely, the ion mass of interest. The first part of the accelerator mass spectrometer is therefore very similar to a conventional mass spectrometer. In this case, the 12C, 13C, and 14C isotopes are analyzed sequentially by varying the voltage on the magnet vacuum box for each isotope in order to analyze each isotope without changing the magnetic field. See ION SOURCES. In the second stage, there is further acceleration by a tandem accelerator. This first accelerates negative ions to the central high-voltage electrode, converts them into positive ions by several successive collisions with gas molecules in a region of higher gas pressure, known as a stripping canal, and then further accelerates the multiply charged positive ions through the same voltage difference back to ground potential. If the acceleration voltage used here is 2 × 106 V, the ions emerging from the accelerator have 8 × 106 eV in energy when they are triply charged in the positive-ion stage. In the third stage, the accelerated ions are analyzed further by the high-energy mass spectrometer. Following an electrostatic analyzer and a magnetic analyzer, the ion currents of 12C and 13C ions can be measured in the two Faraday cups; additional electrostatic or magnetic analysis then completes the isolation of the rare ionic species. In a typical arrangement (see illus.), the currents of the isotopes 12C and 13C are measured in sequence for a short time, and the much rarer 14C ions then are counted for a longer time in a suitable ion detector. The 14C/12C and 13C/12C ratios can then be determined from these measurements. Finally, these ratios determine the age of the sample. Distinguishing features. These three stages of analysis illustrate the features that clearly distinguish accelerator mass spectrometry from conventional mass spectrometry. These are the elimination of molecular ions and isobars from the mass spectrometry. Elimination of molecular ions. The abundant molecular ions, such as 13CH− and 12CH2−, which have similar masses to the 14C− ions, either must be resolved by a very high-resolution mass spectrometer or must be separated by some other method. A tandem accelerator provides a convenient way of completely eliminating molecular ions from the mass spectrometry because ions of a few megaelectronvolts can lose several electrons on passing through the region of higher gas pressure in the stripping canal. Molecules with more than two electrons missing have not been observed in this case, so that the accelerator mass spectrometry of radiocarbon, utilizing charge-3 ions, is free of molecular interferences. The molecular fragments are removed by the high-energy mass spectrometry. Elimination of isobars. The use of a negative-ion source, which is necessary for tandem acceleration, can also
ensure the complete separation of atoms of nearly identical mass (isobars). In the case of radiocarbon analysis, the abundant stable 14N ions and the very rare radioactive 14C ions, which differ in mass by only 1 part in 86,000, are separated completely because the negative ion of nitrogen is unstable whereas the negative ion of carbon is stable. Other examples of rare radioisotopes that can be completely separated at the negative-ion source are 26Al and 129I. In addition, it is frequently possible to count the ions without any background in the ion detectors because of their high energy. In many cases, it is also possible to identify the ion. For example, ions of the isobars 10Be and 10B can easily be distinguished by their differing ranges in matter in the final ion detection system. Both beryllium and boron form negative ions so they cannot be distinguished at the ion source. Advantages. As a result of the destruction of molecules in the accelerator and the elimination of the interfering isobars at the negative-ion source or their separation by the ion detection system, the modest mass-resolution requirements of accelerator mass spectrometry can be achieved with high efficiency. This is a great advantage of accelerator mass spectrometry in the measurement of rare radioactive species. For the study of many rare radioactive atoms, accelerator mass spectrometry also has the important advantage that there can be no background except for contamination with the species being studied. For example, significant interference with the betaparticle counting of radiocarbon from cosmic rays and natural radioactivity occurs for carbon samples about 25,000 years old. In contrast, accelerator mass spectrometer measurements are affected only by the natural contamination of the sample which becomes serious for samples about 50,000 years old. Applications. Accelerator mass spectrometry has been applied to an increasing variety of problems. In particular, radiocarbon dating is carried out at many laboratories, with thousands of dates being determined each year. These dates approach the accuracy of those produced by the beta-particle counting method. The principal advantage of accelerator mass spectrometry for radiocarbon dating is the ability to date smaller samples of carbon and to do it more quickly. The advantage of the lower background is very useful, and the background limit has been the subject of further studies. An example of the use of accelerator mass spectrometry can be found in the dating of a small piece of wood, cut by a metal implement, from an archeological site in Newfoundland. This confirmed directly that the site had been occupied by people from Europe about the year 1000. It was presumably one of the early Viking settlements. There has been an international study of an artifact known as the Shroud of Turin. Tradition holds that this was the burial cloth of Jesus; however, it is now widely believed to be an icon dating from about 1300. Only very small samples from this valuable artifact were made available for dating. See ART CONSERVATION CHEMISTRY.
51
52
Accelerometer Radiocarbon dating has had an increasing impact on geophysical measurements and is particularly suitable for studying the events associated with the most recent ice age. For example, it is possible to date raised beaches by measuring the radiocarbon from single shells instead of being forced to use heterogeneous mixtures of a large number of shells. Accelerator mass spectrometry is used for oceanographic studies, such as measuring the age of the dissolved carbon dioxide (CO2) in ocean water, and the age of foraminifera from ocean sediments. Much smaller samples of water or foraminifera can be collected for analysis than would otherwise be possible, facilitating important studies—for example, of global ocean circulation—that are closely related to the understanding of climate. Radiocarbon has been used extensively as a biomedical tracer in order to detect the effects of toxins, mutagens, carcinogens, or chemotherapeutics in living systems. Radiocarbon is used because it can readily be incorporated into molecules without causing a change in chemical behavior. Accelerator mass spectrometry is important for these studies because of the small sample sizes that can be used, which result in a much lower radiation dose being received by the organism. Beryllium-10. The long-lived radioactive beryllium isotope 10Be is created in the atmosphere by cosmic rays. Since its half-life is 1.6 × 106 years, it is extremely difficult to study by beta-particle counting, but it was first studied at natural abundances by accelerator mass spectrometry, using a cyclotron as an accelerator and range separation to eliminate the abundant natural isobar 10B. It is now measured routinely by accelerator mass spectrometry in many laboratories by using tandem accelerators. The 10Be atoms in sediments subducted beneath continents and later emitted from volcanoes have been used in studies of the time taken for such subterranean transport. Studies of 10Be in ice cores and sediments are used in determining accumulation rates over millions of years. Chlorine-36. The long-lived isotope of chlorine 36Cl with a half-life of 350,000 years has also been extensively studied by accelerator mass spectrometry. Unlike 10Be, the 36Cl atoms readily dissolve in water and so are of importance to the science of hydrology. The 36Cl is also created by cosmic rays in the atmosphere and so can be used to study the transport of water from the atmosphere and through underground aquifers. The 36Cl, like 26Al, is also created in surface rocks by cosmic rays, making possible the dating of the exposure of glacial moraines and certain types of rock. For example, the rocks ejected from a meteor crater have been used to estimate that the crater was formed about 30,000 years ago. Iodine-129. The isotope 129I has a half-life of 1.6 × 107 years and can be studied most easily by accelerator mass spectrometry. The 129I is created in the atmosphere by cosmic rays but is generated mainly by the neutron-induced fission of the uranium isotope 235U. This fission occurs naturally in uranium ores, during the production of nuclear power, and from the fallout
following the use of nuclear weapons. This isotope has been used to study the dispersal in the oceans of radioactive wastes from nuclear fuel reprocessing sites and reactors in Europe and North America. The 129 I has been measured over a wide area of the North Atlantic, and it can be detected easily in liter quantities of water. It has also been studied in connection with the dispersal of radioactive wastes in Russia and the seas north of Russia. Many applications of accelerator mass spectrometry will continue to be made in the future. The small sample size required and the high sensitivity for detecting certain long-lived radioactive atoms or very rare stable atoms such as iridium will facilitate their use as environmental and biological tracers. Methods for the detection of very rare atoms in highpurity materials are also being developed. See DATING METHODS; GEOCHRONOMETRY; MASS SPECTROMAlbert E. Litherland ETRY. Bibliography. L. Brown, Applications of accelerator mass spectrometry, Annu. Rev. Earth Planet. Sci., 12:39–59, 1984; C. E. Chen et al. (eds.), Accelerator Mass Spectrometry, 1992; D. Elmore and F. M. Phillips, Accelerator mass spectrometry for measurement of long-lived radioisotopes, Science, 236:543– 550, 1987; A. E. Litherland, Fundamentals of accelerator mass spectrometry, Phil. Trans. Roy. Soc. London, A323:5–21, 1987.
Accelerometer A mechanical or electromechanical instrument that measures acceleration. The two general types of accelerometers measure either the components of translational acceleration or angular acceleration. See ACCELERATION MEASUREMENT. Translational accelerometers. Most translational accelerometers fall into the category of seismic instruments, which means the accelerations are not measured with respect to a reference point. Of the two types of seismic instruments, one measures the attainment of a predefined acceleration level, and the other measures acceleration continuously. Figure 1 shows the first type of instrument, in which a seismic mass is suspended from a bar made of brittle material which fails in tension at a predetermined
brittle tension member
direction of acceleration mass
Fig. 1. Brittle-member acceleration-level indicator.
Accelerometer output
input
spring
mass direction of acceleration damper
slide-wire
Fig. 2. Slide-wire potentiometer.
seismic mass
direction of acceleration damper
pretensioned wire Fig. 3. Unbounded strain gage accelerometer.
acceleration level. See SEISMOGRAPHIC INSTRUMENTATION. Continuously measuring seismic instruments are composed of a damped or an undamped springsupported seismic mass which is mounted by means of the spring to a housing. The seismic mass is restrained to move along a predefined axis. Also provided is some type of sensing device to measure acceleration. To assure that the deflection of the seismic mass is proportional to the acceleration, such accelerometers are rated so that the excitation frequency associated with the acceleration is only 0.1 to 0.4 of the instrument’s resonance frequency, depending on the amount of damping provided. For most accelerometers the damping coefficient varies between 0 and 0.7 of critical damping. The type of sensing device used to measure the acceleration determines whether the accelerometer is a mechanical or an electromechanical instrument. One type of mechanical accelerometer consists of a liquid-damped cantilever spring-mass system, a shaft attached to the mass, and a small mirror mounted on the shaft. A light beam reflected by the mirror passes through a slit, and its motion is recorded on moving photographic paper. The type of electromechanical sensing device classifies the accelerometer as variable-resistance, variable-inductance, piezoelectric, piezotransistor, or servo type of instrument or transducer. The sensing device for variable-resistance accelerometers operates on the principle that electrical resistance of an conductor is a function of its
dimensions. When the dimensions of the conductor are varied mechanically, as constant current flows through it, the voltage developed across it varies as a function of this mechanical excitation. One type of variable-resistance transducer makes use of a slidewire potentiometer (Fig. 2). This type of accelerometer has low acceleration-measuring ability and low resonance frequency. A more common type of variable-resistance transducer makes use of the fact that the resistance in a wire (called a strain gage) changes when it is strained. Most strain gage accelerometers are either of the unbounded or bounded type. In unbounded strain gage instruments fine pretensioned wires support a fluid-damped seismic mass, and the wires act as both the spring element and the strain gage (Fig. 3). The bounded strain gage transducers have the strain gages mounted directly on the spring element, and the gages measure the strain that takes place in the spring due to the relative deflection of the seismic mass. Wire strain gage accelerometers are characterized by limited high resonance frequency, acceleration measurement up to 1000 g, and ability to measure acceleration signals down to zero (dc) frequencies. See POTENTIOMETER. Another type of variable-resistance accelerometer uses as the spring a piezoresistive element which is a crystal semiconductor whose resistivity changes with applied force (Fig. 4). This type of instrument has much greater sensitivity and resonance frequency than the wire strain gage transducer and is capable of measuring accelerations of 10,000 g or less. Variable-resistance instruments are usually arranged electrically to form a Wheatstone bridge. The bridge input excitation is from dry cells or commercial power supplies, and the output which is usually amplified is recorded by means of an oscilloscope or an oscillograph recorder. A third type of variable-resistance accelerometer employs the Rolamite concept. This accelerometer is composed of two rolling cylinders supported by a thin flexible band, and this low-friction mechanical system is constrained by two parallel guide surfaces (Fig. 5). The flexible band has a cutout which provides a negative spring against the displacement of the rollers, which are also the seismic mass. The sensing device consists of a roller making electrical contact across two resistive elements located on one of
clamp
seismic mass direction of acceleration
piezoresistive or piezoelectric disk Fig. 4. Piezoresistive or piezoelectric accelerometer.
53
54
Acceptor atom roller
flexible band
resistive element
band cutout guide surface
direction of acceleration
diaphragm seismic mass
Fig. 5. Rolamite type of accelerometer.
differential transformer
turn, cause a very large reversible change in current across the underlying pn junction of the transistor. This change of current is a measure of the acceleration. These instruments are characterized by small size, very high resonance frequencies, acceleration measurements of 400 g or less, and voltage outputs that are high enough to eliminate need for amplifiers.
electromagnetic damping
output stylus
direction of acceleration
transistor
Fig. 7. Piezotransistor accelerometer.
input
direction of acceleration seismic mass
spring
Fig. 6. Variable-inductance accelerometer.
the guide surfaces. The Rolamite accelerometer can be designed to be miniature in size and to measure accelerations in the range 1–500 g. The variable-inductance accelerometer is a differential transformer which makes use of three coils arranged as shown in Fig. 6. The center coil is excited from an external ac power source, and the two end coils connected in series opposition are used to produce an ac output which is proportional to the displacement of a seismic mass passing through the coils. This instrument is characterized by fairly low resonance frequencies, acceleration measurements of 50 g or less, and high voltage outputs at low frequencies which eliminate the need for preamplifiers. Piezoelectric accelerometers utilize a piezoelectric crystal which supports the seismic mass in such a way that the crystal acting as a spring is strained in either compression, flexure, or shear (Fig. 4). Many crystals, such as quartz, Rochelle salts, and barium titanate, have these properties. A number of ceramic materials can be used when they contain certain impurities. Piezoelectric instruments are noted for their small size, large acceleration measurement up to 5000 g, very high resonance frequency, inability to measure dc frequency acceleration signals, and the need for specialized circuitry. The acceleration signal is usually amplified by means of a charge or voltage amplifier and then recorded. A piezotransistor accelerometer consists of a seismic mass supported by a stylus which transmits a concentrated force to the upper diode surface of the transistor (Fig. 7). Because of this force, localized stresses occur in the upper diode surface which, in
A servo accelerometer is a high-accuracy transducer, containing an electromechanical servomechanism which holds a mass in a fixed position as the instrument is accelerated. The force required to restrain the mass is a measure of the acceleration. This instrument has limited high-frequency response, acceleration outputs of 100 g or less, and capability of measuring dc frequency responses. See SERVOMECHANISM. Angular accelerometers. There are several different types of angular accelerometers. In one type the damping fluid serves as the seismic mass. Under angular acceleration the fluid rotates relative to the housing and causes on two symmetrical vanes a pressure which is a measure of the angular acceleration. Another type of instrument has a fluid-damped symmetrical seismic mass in the form of a disk which is so mounted that it rotates about the normal axis through its center of gravity. The angular deflection of the disk, which is restrained by a spring, is proportional to the angular acceleration. The sensing devices used in these instruments to measure the angular acceleration are of the variable-resistance and servo types. Roger C. Duffield; Teruo Ishihara Bibliography. T. G. Beckwith, R. D. Marangoni, and J. H. Lienhard, Mechanical Measurements, 6th ed., 2006; W. Boyes, Instrumentation Reference Book, 3d ed., 2002; E. O. Doebelin, Measurement Systems: Application and Design, 5th ed., 2003; C. P. Wright, Applied Measurement Engineering, 1994.
Acceptor atom An impurity atom in a semiconductor which can accept or take up one or more electrons from the crystal and become negatively charged. An atom which substitutes for a regular atom of the material but has one less valence electron may be expected to be an acceptor atom. For example, atoms of boron, aluminum, gallium, or indium are acceptors in germanium and silicon (illus. a), and atoms of antimony
Acetal energy
Si
conduction band
Eg
B− Si Si
EI
0
+
valence band
Si (b)
(a)
Trivalent acceptor atom, boron (B), in the elemental semiconductor silicon (Si). (a) Boron atom in a substitutional position, that is, replacing silicon, a tetravalent host atom, by completing the four tetrahedral covalent bonds with its nearest neighbor silicon atoms. This requires an electron to be accepted from the valence band, thus making boron negatively charged. (b) Energy diagram showing that the absence of an electron in the valence band is equivalent to a positive charge carrier, a hole, which is bound to boron via Coulomb attraction with an ionization enerzy EI. Eg = energy gap separating valence band from conduction band.
and bismuth are acceptors in tellurium crystals. Acceptor atoms tend to increase the number of holes (positive charge carriers) in the semiconductor (illus. b). The energy gained when an electron is taken up by an acceptor atom from the valence band of the crystal is the ionization energy of the atom. See DONOR ATOM; SEMICONDUCTOR. H. Y. Fan; Anant K. Ramdas Bibliography. C. Kittel, Introduction to Solid State Physics, 7th ed., Wiley, 1996; A. K. Ramdas and S. Rodriguez, Spectroscopy of the solid-state analogues of the hydrogen atom: Donors and acceptors in semiconductors, Rep. Prog. Phys., 44:1297–1387, 1981.
Acetal A geminal diether (R1 = H). Ketals, considered a subclass of acetals, are also geminal diethers (R1 = C, aliphatic or aromatic). Acetals are (1) independent structural units or a part of certain biological and commercial polymers, (2) blocking or protecting groups for complex molecules undergoing selective synthetic transformations, and (3) entry compounds for independent organic chemical reactions. Acetals are easily prepared by the reaction of aldehydes with excess alcohol, under acid-catalyzed conditions. This is usually a two-step process, as shown below, in which an aldehyde is treated with an OH
OR2
R2OH
R
C
R2OH
R
O
C
OR2
H2O
R
C
OR2
H2O
R1
R1
R1
Aldehyde or ketone
Hemiacetal or hemiketal
Acetal or ketal
alcohol to yield a less stable hemiacetal, which then
reacts with additional alcohol to give the acetal. Protonic or Lewis acids are effective catalysts for acetal formation; dehydrating agents, such as calcium chloride and molecular sieves, can also be used for molecules, such as sugars, where acids may cause problems. Less common acetal preparations are Grignard reagent condensation with orthoformates and mercuric-catalyzed additions of alcohols to acetylenes. See GRIGNARD REACTION. Two well-known hemiacetals (monomers) are α- and β-glucose: the 1,4-polysaccharide from αglucose is amylose (amylopectin has an acetal 1,6linkage branch from the linear polymer backbone), and the 1,4-polysaccharide from β-glucose is cellulose. The acetal link connects the monomers incorporated into the polymer in each case, and it is formed by the enzymatic condensation of the hemiacetal functional group (C1 OH) of one monomer with the (C4 OH) hydroxyl functional group of another monomer. Several other biopolymers containing similar acetal links between repeating units are chitin, inulin, and pectin. Synthetic polymers with 1,3-diol units, such as polyvinyl alcohol, can be transformed into polymeric acetals, which are stable and useful plastics (for example, formal and butyral). Difunctional aldehydes and ketones can be condensed with certain tetrols (for example, pentaerythritol) for preparation of polyketals, and polyspiroketals, which are highly crystalline materials. Polyoxymethylene, which is the acetal resin homopolymer of formaldehyde, is prepared by cationic polymerization of trioxane, the cyclic acetal of formaldehyde. Homopolymerization of acetaldehyde or acetone, where the acetal or ketal oxygen is also part of the polymer backbone, is possible, but requires specialized conditions. See CELLULOSE; POLYACETAL; POLYMER; POLYMERIZATION; POLYSACCHARIDE; STARCH. Their use as protective groups is attributed to the fact that acetals (or ketals) are readily cleaved by aqueous acids to the respective aldehyde (or ketone) and that they are stable to alkalies and nucleophiles. Cyclic acetals, such as meta-dioxolanes and meta-dioxanes, are very convenient to work with for blocking–unblocking types of syntheses. During these multistep processes, the carbonyl portion of the molecule is transformed to its acetal (or ketal), while synthetic reactions that do not utilize acid are effected on other functional groups of the molecule. The acetal (or ketal) part of the molecule is then readily hydrolyzed to free the carbonyl functional group. The cyclic acetals can also be used as aprotic solvents and plasticizers. See ALDEHYDE; CARBONYL; KETONE; ORGANIC SYNTHESIS. Acetals undergo a variety of useful reactions such as α-acylation, α-bromination, condensations with enol acetates and silyl enolates, conversion to enol ethers, oxidation to esters, reduction to ethers, and reactions with active hydrogen compounds (for example, diethyl malonate). Acetal additions to vinyl ethers followed by hydrolysis result in a preparation of α,β-unsaturated aldehydes. Dioxolanes (from aldehydes) are also treated with alkyllithiums for the preparation of ketones.
55
56
Acetic acid Thioacetals, especially cyclic thioacetals, can be reduced to alkanes with Raney nickel, and 1,3dithianes, prepared from aldehydes, can be metalated, condensed with a variety of electrophilic reagents, and hydrolyzed to free the functionalized aldehyde. Charles F. Beam Bibliography. R. J. Fessenden and J. S. Fessenden, Organic Chemistry, 6th ed., 1998; D. S. Kemp and F. Vellaccio, Organic Chemistry, 1980; R. W. Lenz, Organic Chemistry of Synthetic High Polymers, 1967, reprint 1988; J. March, Advanced Organic Chemistry: Reactions, Mechanisms, and Structures, 4th ed., 1992; R. B. Seymour and C. E. Carraher, Jr., Polymer Chemistry, 1982; M. P. Stevens, Polymer Chemistry, 3d ed., 1998.
Acetic acid A colorless, pungent liquid, CH3COOH, melting at 16.7◦C (62.1◦F) and boiling at 118.0◦C (244.4◦F). Acetic acid is the sour principle in vinegar. Concentrated acid is called glacial acetic acid because of its readiness to crystallize at cool temperatures. Acetic acid in vinegar arises through an aerobic fermentation of dilute ethanol solutions, such as wine, cider, and beer, with any of several varieties of Acetobacter. Chemical properties. Pure acetic acid is completely miscible with water, ethanol, diethyl ether, and carbon tetrachloride, but is not soluble in carbon disulfide. Freezing of acetic acid is accompanied by a remarkable volume contraction: the molar volume of liquid acetic acid at the freezing point is 57.02 cm3/mole, but at the same temperature the crystalline solid is 47.44 cm3/mole. It is a strongly proton-donating solvent with a relatively small dipole moment and a low dielectric constant. In a water solution, acetic acid is a typical weakly ionized acid (Ka = 1.8 × 10−5). See CARBOXYLIC ACID. The vapor density of acetic acid indicates a molecular weight considerably higher than would be expected for a compound with a formula weight of 60.05. The acid probably exists largely as the dimer in the vapor and liquid states. Acetic acid neutralizes many oxides and hydroxides, and decomposes carbonates to furnish acetate salts, which are used in textile dyeing and finishing, as pigments, and as pesticides; examples are verdigris, white lead, and paris green. See ARSENIC; LEAD. Over two-thirds of the acetic acid manufactured is used in production of either vinyl acetate or cellulose acetate. Acetic anhydride, the key intermediate in making cellulose acetate, is prepared commercially by pyrolysis of acetic acid in the presence of trialkyl phosphate catalyst. Considerable amounts of acetic acid are consumed in the manufacture of terephthalic acid by liquid-phase oxidation of xylene, and in the preparation of esters for lacquer solvents, paints and varnishes, pharmaceuticals, and herbicides. Production. Acetic acid was formerly manufactured from pyroligneous acid obtained in destructive distillation of wood. These processes are of historical
interest because many modern chemical engineering operations developed through the study of acetic acid production. Today acetic acid is manufactured by three main routes: butane liquid-phase catalytic oxidation in acetic acid solvent, palladium–copper salt–catalyzed oxidation of ethylene in aqueous solution, and methanol carbonylation in the presence of rhodium catalyst. Large quantities of acetic acid are recovered in manufacturing cellulose acetate and polyvinyl alcohol. Some acetic acid is produced in the oxidation of higher olefins, aromatic hydrocarbons, ketones, and alcohols. See ESTER; OXIDATION PROCESS; SOLVENT; WOOD CHEMICALS. Frank Wagner Bibliography. W. H. Brown, Introduction to Organic and Biochemistry, 4th ed., 1987; T. A. Geissman, Principles of Organic Chemistry, 4th ed., 1977; J. D. Roberts and M. C. Caserio, Basic Principles of Organic Chemistry, 2d ed., 1977.
Acetone A chemical compound, CH3COCH3; the first member of the homologous series of aliphatic ketones. It is a colorless liquid with an ethereal odor. Its physical properties include boiling point 56.2◦C (133◦F), melting point −94.8◦C (−138.6◦F), and specific gravity 0.791. Acetone is an extremely important, lowcost raw material that is used for production of other chemicals. Acetone is used as a solvent for cellulose ethers, cellulose acetate, cellulose nitrate, and other cellulose esters. Cellulose acetate is spun from acetone solution. Lacquers, based on cellulose esters, are used in solution in mixed solvents including acetone. Acetylene is safely stored in cylinders under pressure by dissolving it in acetone, which is absorbed on inert material such as asbestos. It has a low toxicity. Production. The principal method of acetone production uses propylene, obtained from the cracking of petroleum. Addition of sulfuric acid to propylene yields isopropyl hydrogen sulfate, which upon hydrolysis yields isopropyl alcohol. Oxidation or dehydrogenation over metal catalysts, such as copper, converts the alcohol to acetone, as shown in reaction (1). CH2
CH3 H2O
H2SO4
CH
CHOSO2OH
CH3
CH3 CH3
CH3 Cu
CHOH
C
O
(1)
heat
CH3
CH3
Acetone is also produced by passage of acetic acid vapor over metallic oxide catalysts at 400–450◦C (750–840◦F), by partial oxidation of the lower alkane hydrocarbons, and by the decomposition of cumene
Acetylcholine hydroperoxide, as shown in reaction (2). Phenol is the other product of this last process. OOH CH3CHCH3
CH3
C
CH3 O2
heat acid
that the perfusate from a frog’s heart that was slowed by stimulation of the vagus nerve slowed the heart of a second frog. The active substance in the perfusate was expected to be acetylcholine, which was known to slow the heart; this proved to be correct. Biosynthesis. Acetylcholine is synthesized at axon endings from acetyl coenzyme A and choline by the enzyme choline acetyltransferase which transfers the acetyl group to choline, reaction (1). Acetylcholine Acetyl coenzyme A + choline−−−−→
OH
coenzyme A + acetylcholine
+ CH3COCH3
(2)
Chemical uses. Pyrolysis of acetone vapor at 700◦C (1300◦F) produces ketene, which reacts with acetic acid to produce acetic anhydride, as shown in reaction (3).
CH2
is stored at each ending in hundreds of thousands of membrane-enclosed synaptic vesicles. The axon ending is also rich in mitochondria that supply acetyl coenzyme A using pyruvate (from glucose) as the source of the acetyl group, in a reaction catalyzed by the pyruvate dehydrogenase complex, reaction (2). Pyruvate + coenzyme A + NAD+ −−−−→
heat
CH3COCH3
(1)
C
O
CH3COOH
(CH3CO)2O
(3)
Aldol-type condensation of acetone with Ba(OH)2 yields diacetone alcohol (1), mesityl oxide (2), and phorone (3), shown in reaction (4). Catalytic hydro-
acetyl coenzyme A + CO2 + NADH, H+
(2)
Mitochondria also supply adenosine triphosphate for the acetylation of coenzyme A by acetate in a twostep reaction (3). Acetate + adenosine triphosphate
Ba(OH)2
2CH3COCH3
acetyl adenosine monophosphate + pyrophosphate CH3
OH
+ coenzyme A
H+
CH3CCH2COCH3
CH3C
CHCOCH3
adenosine monophosphate + acetyl coenzyme A
(2)
CH3 (1)
CH3 CH3C
CH3 CHCOCH
CCH3
(4)
(3)
genation of mesityl oxide gives methyl isobutyl ketone. All these products are of commercial use as solvents or as chemical intermediates. See KETENE; KETONE. David A. Shirley
Acetylcholine A naturally occurring quaternary ammonium cation ester, with the formula CH3(O)COC2H4N(CH3)3+, that plays a prominent role in nervous system function. Neurotransmission. The great importance of acetylcholine derives from its role in physiology as a neurotransmitter for cholinergic neurons. Cholinergic nerves innervate many tissues, including smooth muscle and skeletal muscle, the heart, ganglia, and glands. The effect of stimulating a cholinergic nerve, for example, the contraction of skeletal muscle or the slowing of the heartbeat, results from the release of acetylcholine from the nerve endings. Other types of nerves release other transmitters, for example, norepinephrine. The release of a neurotransmitter was first clearly shown by Otto Loewi, who demonstrated in 1921
(3)
Molecular biology of axon. Neurons, like other cells, have a high internal concentration of potassium ions and a low internal concentration of sodium ions maintained by a sodium–potassium “pump.” An axon at rest is permeable mainly to potassium, and so a very small amount of potassium leaves the cell and establishes a membrane potential of about –80 millivolts (the inside negative) that nearly blocks further leakage of potassium. This small potential produces a strong electric field in the membrane because the membrane is thin and has a low dielectric constant. The potential and high external concentration of sodium ions could drive a large inward sodium current, if the membrane were permeable to sodium. See ION TRANSPORT. Axonal membranes contain voltage-gated sodium and potassium channels. These channels are composed of proteins that span the membrane and allow the transient flow of their specific ion, depending on the membrane potential. When an axon is stimulated at some point by slightly increasing the membrane potential to a less negative value, there is a transient opening of sodium channels followed by a transient opening of potassium channels. As a result, sodium ions enter the cell and drive the membrane potential positive. As the sodium current subsides, the potassium channels open and the outward flow of potassium restores the resting potential. This sequence of potential changes, the action potential, is self-propagating as a nerve impulse that moves along
57
58
Acetylcholine the axon away from the cell body toward the numerous branched endings (toward the cell body in a dendrite). Thus, although the electric currents are perpendicular to the axon, the impulse moves parallel to the axon. See BIOPOTENTIALS AND IONIC CURRENTS. Molecular biology of synapse. When an impulse reaches an axon ending, voltage-gated calcium channels open and calcium, which is extremely low inside the cell, enters the nerve ending. The increase in calcium-ion concentration causes hundreds of synaptic vesicles to fuse with the cell membrane and expel acetylcholine into the synaptic cleft (exocytosis). The acetylcholine released at a neuromuscular junction binds reversibly to acetylcholine receptors in the muscle end-plate membrane, a postsynaptic membrane that is separated from the nerve ending by a very short distance. The receptor is a cation channel. In Torpedo, an electric fish, the receptor is composed of five protein subunits of four kinds, α 2βγ δ. The channel opens when two acetylcholine molecules are bound, one to each α subunit. A sodium current enters the cell and depolarizes the membrane, producing an end-plate potential which triggers an action potential. The action potential, as in the nerve, is self-propagating. The resulting impulse indirectly causes the muscle to contract. Acetylcholine must be rapidly removed from a synapse in order to restore it to its resting state. This is accomplished in part by diffusion but mainly by the enzyme acetylcholinesterase which hydrolyzes acetylcholine, reaction (4). The enzyme mechanism Acetylcholine + H2 O−−−−→acetic acid + choline
(4)
involves the transfer of the acetyl group from choline to the hydroxyl group of a specific serine residue of the enzyme located at the active site. The reaction is completed when the acetyl enzyme, itself an ester, hydrolyzes to form acetic acid and free enzyme. As expected, acetylcholinesterase is a very fast enzyme: one enzyme molecule can hydrolyze 10,000 molecules of acetylcholine in 1 s. Red cells contain acetylcholinesterase and blood serum contains a different cholinesterase, but the function of these enzymes is not clearly known. See SYNAPTIC TRANSMISSION. Poisons, insecticides, and drugs. It is evident that any substance that efficiently inhibits acetylcholinesterase will be extremely toxic. Thus, synthetic organophosphonates such as Sarin, Tabun, and Soman are military nerve gases. These substances, which are volatile liquids, transfer a phosphonyl group to the specific serine side chain of the enzyme to produce a phosphonyl-enzyme derivative analogous to the acetyl enzyme. But unlike the acetyl enzyme, the phosphonyl enzyme hydrolyzes only extremely slowly and is trapped as an inactive derivative. Some organophosphates (similar to phosphonates) are useful: Parathion and Malathion are selectively toxic to insects and are important insecticides. Organophosphate poisoning is treated with atropine and 2-pyridine aldoxime methiodide
(2-PAM): atropine antagonizes some of the effects of acetylcholine, and 2-PAM reactivates the inhibited enzyme by removing the phosphoryl group. See INSECTICIDE. Similarly, certain carbamates form carbamyl derivatives of acetylcholinesterase. Carbaryl is an important carbamate insecticide. An increase in the persistence of acetylcholine at synapses is sometimes medically desirable, and carbamates are used to treat a number of medical disorders. Pyridostigmine is extensively used in the treatment of myasthenia gravis, a disease characterized by muscle weakness which is believed to be an autoimune disease, that is, antibodies are produced against the acetylcholine receptor. The proper dosage of the carbamate increases the persistence of acetylcholine at the neuromuscular junction and improves muscle strength. Carbamates, once extensively used in treating glaucoma, usually in conjunction with pilocarpine which activates the acetylcholine receptor now find limited use. Other types of acetylcholinesterase inhibitors are used in managing Alzheimer’s disease. See MYASTHENIA GRAVIS. Succinyldicholine, which reacts with the acetylcholine receptor of skeletal muscle, is used to produce paralysis and relaxation of muscle during surgery. This ester is not hydrolyzed by acetylcholinesterase but is hydrolyzed by serum cholinesterase which limits the duration of paralysis. Because a small fraction of the population has a variant cholinesterase, patients who will be given succinyldicholine are tested to be sure their serum will hydrolyze this drug. There are a number of exotic, naturally occurring poisons that interfere with the nerve impulse or with synaptic transmission. Many of these substances have been very useful in studying the mechanism of these processes. Botulinus toxin prevents the release of acetylcholine, and black widow spider toxin causes the release of acetylcholine. The jungle-dart poison curare binds to the acetylcholine receptor of skeletal muscle and causes paralysis. Bungarotoxin and cobratoxin found in snake venoms have a similar effect. Muscarine from some poisonous mushrooms and arecoline from the betel nut react with the acetylcholine receptor of other tissues. Tetrodotoxin from the Japanese puffer fish, saxitoxin from the bacterium that causes the “red tide” of mussels and clams, and batrachotoxin from a South American frog bind to and block the function of sodium channels, thus preventing the nerve impulse. See POISON; TOXICOLOGY; TOXIN. Electric organs. Some aquatic animals stun or kill their prey by producing a high-voltage electrical discharge. Their electric organs are composed of thousands of flat cells arranged in series. These cells are derived from skeletal muscle and have a conducting membrane on one face which also contains many nerve junctions. The nearly simultaneous excitation of these faces produces action potentials that are summed by the series arrangement to produce hundreds of volts. See ELECTRIC ORGAN (BIOLOGY); NERVOUS SYSTEM (VERTEBRATE). Irwin B. Wilson
Acid and base Bibliography. B. Alberts et al., Molecular Biology of the Cell, 1994; A. G. Gilman et al. (eds.), The Pharmacological Basis of Therapeutics, 9th ed., 1996; L. Stryer, Biochemistry, 4th ed., 1995.
HC
CCH2OH CH2O
CHCl
Vinyl chloride
CCH2OH Butynediol
2CH2O H2O
HCI
H2C
Acetylene
HOCH2C
HC
CH
CH3CHO Acetaldehyde O
CO H2O
An organic compound with the formula C2H2 or HC CH. The first member of the alkynes, acetylene is a gas with a narrow liquid range; the triple point is −81◦C (−114◦F). The heat of formation (H◦f) is +227 kilojoules/mole, and acetylene is the most endothermic compound per carbon of any hydrocarbon. The compound is thus extremely energyrich and can decompose with explosive force. At one time acetylene was a basic compound for much chemical manufacturing. It is highly reactive and is a versatile source of other reactive compounds. Preparation. The availability of acetylene does not depend on petroleum liquids, since it can be prepared by hydrolysis of calcium carbide (CaC2), obtained from lime (CaO), and charcoal or coke (C) [reaction (1)]. Other preparative methods are based H2 O
CaO + 3C −→ CO + CaC2 −−−−−→ C2 H2 + Ca(OH)2
(1)
on pyrolysis of some carbon source. In modern practice, methane (CH4) is passed through a zone heated to 1500◦C (2732◦F) for a few milliseconds, and acetylene is then separated from the hydrogen in the effluent gas [reaction (2)]. 2CH4 −→ C2 H2 + 3H2
CH3COH
Ni
H2C
CHCO2H
H 2C CuCl
Acrylic acid
O
Vinyl acetate
HCN
Cl HCl
H2C
CHCN
Acrylonitrile
HC
C
CH
CH2
Vinyl acetylene
H2C
CHC Chloroprene
Because of the very high heat of formation and combustion, an acetylene-oxygen mixture provides a very high temperature flame for welding and metal cutting. For this purpose acetylene is shipped as a solution in acetone, loaded under pressure in cylinders that contain a noncombustible spongy packing. See ALKYNE; TORCH. James A. Moore Bibliography. Kirk-Othmer Encyclopedia of Chemical Technology, vol. 1, 4th ed., 1991; Ullmann’s Encyclopedia of Industrial Chemistry, vol. A1, 5th ed., 1985.
(2)
Uses. There are processes for the conversion of acetylene to a number of important compounds, including vinyl chloride, vinyl acetate, acrylonitrile, and chloroprene, all of which are monomer intermediates for polymer production [reaction (3)]. Much of this chemistry is applicable also to higher alkynes. Although acetylene is a versatile and flexible raw material, nearly all of these processes became obsolete as the petrochemical industry developed and ethylene became a more economical feedstock for large-volume vinyl monomers and other C2 chemicals. A major consideration is that acetylene is a very dangerous substance and handling costs are high. Production of acetylene dropped about 10% per year from a high of 1.1 × 109 lb in 1969 to 3.5 × 108 lb (1.6 × 108 kg) in 1990. See ETHYLENE. The main use of acetylene is in the manufacture of compounds derived from butyne-1,4-diol. The latter is obtained by condensation of acetylene with two moles of formaldehyde and is converted to butyrolactone, tetrahydrofuran, and pyrrolidone. Two additional products, vinyl fluoride and vinyl ether, are also based on acetylene. Cyclooctatetraene, the tetramer of acetylene, is potentially available in quantity if a demand should develop. A superior grade of carbon black that is obtained by pyrolysis of acetylene above 1500◦C (2732◦F) has good electrical conductivity and is produced for use in dry-cell batteries. See PYROLYSIS.
(3)
CHOCCH3
Acid and base Two interrelated classes of chemical compounds, the precise definitions of which have varied considerably with the development of chemistry. These changing definitions have led to frequent controversies, some of which are still unresolved. Acids initially were defined only by their common properties. They were substances which had a sour taste, which dissolved many metals, and which reacted with alkalies (or bases) to form salts. For a time, following the work of A. L. Lavoisier, it was believed that a common constituent of all acids was the element oxygen, but gradually it became clear that, if there were an essential element, it was hydrogen, not oxygen. In fact, the definition of an acid, formulated by J. von Liebig in 1840, as “a hydrogen-containing substance which will generate hydrogen gas on reaction with metals” proved to be satisfactory for about 50 years. Bases initially were defined as those substances which reacted with acids to form salts (they were the “base” of the salt). The alkalies, soda and potash, were the best-known bases, but it soon became clear that there were other bases, notably ammonia and the amines. Acids and bases are among the most important chemicals of commerce. The inorganic acids are often known as mineral acids, and among the most important are sulfuric, H2SO4; phosphoric, H3PO4; nitric, HNO3; and hydrochloric, HCl (sometimes
CH2
59
60
Acid and base called muriatic). Among the many important organic acids are acetic, CH3COOH, and oxalic, H2C2O4, acids, and phenol, C6H5OH. The important inorganic bases are ammonia, NH3; sodium hydroxide or soda, NaOH; potassium hydroxide, KOH; calcium hydroxide or lime, Ca(OH)2; and sodium carbonate, Na2CO3. There are also many organic bases, mostly derivatives of ammonia. Examples are pyridine, C5H5N, and ethylamine, C2H5NH2. Arrhenius-Ostwald theory. When the concept of ionization of chemical compounds in water solution became established, some considerably different definitions of acids and bases became popular. Acids were defined as substances which ionized in aqueous solution to give hydrogen ions, H+, and bases were substances which reacted to give hydroxide ions, OH−. These definitions are sometimes known as the Arrhenius-Ostwald theory of acids and bases and were proposed separately by S. Arrhenius and W. Ostwald. Their use makes it possible to discuss acid and base equilibria and also the strengths of individual acids and bases. The ionization of an acid in water can be written as Eq. (1). Qualitatively, an HA = H+ + A−
(1)
acid is strong if this reaction goes extensively toward the ionic products and weak if the ionization is only slight. A quantitative treatment of this ionization of dissociation can be given by utilizing the equilibrium expression for the acid, as shown in Eq. (2), where +
−
[H ][A ] = KHA [HA]
(2)
the brackets mean concentration in moles per liter and the constant KHA is called the dissociation constant of the acid. This dissociation constant is a large number for a strong acid and a small number for a weak acid. For example, at 25◦C (77◦F) and with water as the solvent, KHA has the value 1.8 × 10−5 for a typical weak acid, acetic acid (the acid of vinegar), and this value varies only slightly in dilute solutions as a function of concentration. Dissociation constants vary somewhat with temperature. They also change considerably with changes in the solvent, even to the extent that an acid which is fully ionized in water may, in some other less basic solvent, become decidedly weak. Almost all the available data on dissociation constants are for solutions in water, partly because of its ubiquitous character, and partly because it is both a good ionizing medium and a good solvent. Acetic acid has only one ionizable hydrogen and is called monobasic. Some other acids have two or even three ionizable hydrogens and are called polybasic. An example is phosphoric acid, which ionizes in three steps, shown in Eqs. (3), (4), and (5), each with its own dissociation constant. Ionization reaction H3 PO4 = H+ + H2 PO4 − −
+
H2 PO4 = H + HPO4
2−
HPO4 2− = H+ + PO4 3−
KHA , 25 ◦ C (77 ◦ F) 1.5 × 10−3 6.2 × 10
−8
4 × 10−13
(3) (4) (5)
A similar discussion can be given for the ionization of bases in water. However, the concentrations of the species H+ and OH− in a water solution are not independently variable. This is because water itself is both a weak acid and a weak base, ionizing very slightly according to Eq. (6). For pure water, H2 O = H+ + OH−
(6)
the concentration of H+ and OH− are equal. At ordinary temperatures, roughly 2 × 10−7% of the water is present as ions. As a result of this ionization, the ionic concentrations are related through Eq. (7). At [H+ ][OH− ] =K [H2 O]
(7)
25◦C (77◦F) and with concentrations in moles per liter, the product [H+][OH−] is equal to 1 × 10−14. A major consequence of this interdependence is that measurement of the concentration of either H+ or OH− in a water solution permits immediate calculation of the other. This fact led S. P. L. Sørenson in 1909 to propose use of a logarithmic pH scale for the concentration of hydrogen ions in water. Although there are some difficulties in giving an exact definition of pH, it is very nearly correct for dilute solutions in water to write Eq. (8). It then turns out that pH pH = − log[H+ ]
(8)
values of 0–14 cover the range from strongly acidic to strongly basic solutions. The pH of pure water at ordinary temperature is 7. See WATER. For many situations, it is desirable to maintain the H+ and OH− concentration of a water solution at low and constant values. A useful device for this is a mixture of a weak acid and its anion (or of a weak base and its cation). Such a mixture is called a buffer. A typical example is a mixture of sodium acetate and acetic acid. From the treatment just given, it is evident that for this case Eq. (9) can be formed. In the [H+ ] =
[CH3 COOH] × 1.8 × 10−5 [CH3 COO− ]
(9)
equation [CH3COOH] and [CH3COO−] represent the concentrations of acetic acid and acetate ion, respectively. Thus, if the concentrations of acetic acid and acetate ion are both 0.1 mole per liter, the H+ concentration will be 1.8 × 10−5 mole per liter and OH− will be 5.5 × 10−10 mole per liter. The pH of this solution will be about 4.7. Constant acidity is a most important aspect of blood and other life fluids; these invariably contain weak acids and bases to give the necessary buffering action. ¨ Bronsted theory. The Arrhenius, or water, theory of acid and bases has many attractive features, but it has also presented some difficulties. A major difficulty was that solvents other than water can be used for acids and bases and thus need consideration. For many of the solvents of interest, the necessary extensions of the water theory are both obvious and plausible. For example, with liquid ammonia as the solvent, one can define NH4+ as the acid ion and NH2− as the base ion, and the former can be thought
Acid and base of as a hydrogen ion combined with a molecule of the solvent. However, for a hydrogenless (aprotic) solvent, such as liquid sulfur dioxide, the extensions are less obvious. Consideration of such systems has led to some solvent-oriented theories of acids and bases to which the names of E. C. Franklin and A. F. D. Germann often are attached. The essence of these theories is to define acids and bases in terms of what they do to the solvent. Thus, one definition of an acid is that it gives rise to “a cation which is characteristic of the solvent,” for example, SO2+ from sulfur dioxide. These theories have been useful in emphasizing the need to consider nonaqueous systems. However, they have not been widely adopted, at least partly because a powerful and wide-ranging protonic theory of acids and bases was introduced by J. N. Br¨ onsted in 1923 and was rapidly accepted by many other scientists. Somewhat similar ideas were advanced almost simultaneously by T. M. Lowry, and the new theory is occasionally called the Br¨ onsted-Lowry theory. This theory gives a unique role to the hydrogen ion, and there appeared to be justification for this. One justification, of course, was the historically important role already given to hydrogen in defining acids. A rather different justification involved the unique structure of hydrogen ion. It is the only common ion which consists solely of a nucleus, the proton. As a consequence, it is only about 10−14 cm in diameter. All other ordinary ions have peripheral electron clouds and, as a result, are roughly 106 times larger than the proton. The small size of the latter makes it reasonable to postulate that protons are never found in a free state but rather always exist in combination with some base. The Br¨ onsted theory emphasizes this by proposing that all acid-base reactions consist simply of the transfer of a proton from one base to another. The Br¨ onsted definitions of acids and bases are: An acid is a species which can act as a source of protons; a base is a species which can accept protons. Compared to the water theory, this represents only a slight change in the definition of an acid but a considerable extension of the term base. In addition to hydroxide ion, the bases now include a wide variety of uncharged species, such as ammonia and the amines, as well as numerous charged species, such as the anions of weak acids. In fact, every acid can generate a base by loss of a proton. Acids and bases which are related in this way are known as conjugate acid-base pairs, and the table lists several examples. By these definitions, such previously distinct chem-
Conjugate acid-base pairs Strong acids
Weak acids
H2 SO4 HCl H3 O⫹ HSO⫺ 4 HF(aq) CH3 COOH HN4 ⫹ HCO⫺ 3 H2 O C2 H5 OH
HSO4 ⫺ Cl⫺ H2 O SO4 2⫹ F⫺ CH3 COO⫺ NH3 CO3 2⫹ OH⫺ C2 H5 O⫺
Weak bases
Strong bases
ical processes as ionization, hydrolysis, and neutralization become examples of the single class of proton transfer or protolytic reactions. The general reaction is expressed as Eq. (10). This equation can Acid1 + base2 = base1 + acid2
(10)
be considered to be a combination of two conjugate acid-base pairs, and the pairs below can be used to construct a variety of typical acid-base reactions. For example, the ionization of acetic acid in water becomes Eq. (11). Water functions here as a base to CH3 COOH + H2 O = CH3 COO− + H3 O+ a1 b2 b1 a2
(11)
form the species H3O+, the oxonium ion (sometimes called the hydronium ion). However, water can also function as an acid to form the base OH−, and this dual or amphoteric character of water is one reason why so many acid-base reactions occur in it. As the table shows, strengths of acids and bases are not independent. A very strong Br¨ onsted acid implies a very weak conjugate base and vice versa. A qualitative ordering of acid strength or base strength, as above, permits a rough prediction of the extent to which an acid-base reaction will go. The rule is that a strong acid and a strong base will react extensively with each other, whereas a weak acid and a weak base will react together only very slightly. More accurate calculations of acid-base equilibria can be made by using the ordinary formulation of the law of mass action. A point of some importance is that, for ionization in water, the equations reduce to the earlier Arrhenius-Ostwald type. Thus, for the ionization of acetic acid in water Eq. (11) leads to Eq. (12). [H3 O+ ][CH3 COO− ] =K [CH3 COOH][H2 O]
(12)
Remembering that the concentration of water will be almost constant since it is the solvent, this can be written as Eq. (13), where KHAc is just the conven[H3 O+ ][CH3 COO− ] = K[H2 O] ≡ K HAc [CH3 COOH]
(13)
tional dissociation constant for acetic acid in water. One result of the Br¨ onsted definitions is that for a given solvent, such as water, there is only a single scale of acid-base strength. Put another way, the relative strength of a set of acids will be the same for any base. Hence, the ordinary tabulation of ionization constants of acids in water permits quantitative calculation for a very large number of acid-base equilibria. The Br¨ onsted concepts can be applied without difficulty to other solvents which are amphoteric in the same sense as water, and data are available for many nonaqueous solvents, such as methyl alcohol, formic acid, and liquid ammonia. An important practical point is that relative acid (or base) strength turns out to be very nearly the same in these other solvents as it is in water. Br¨ onsted acid-base reactions can also be studied in aprotic solvents (materials such as
61
62
Acid and base hexane or carbon tetrachloride which have virtually no tendency to gain or lose protons), but in this case, both the acid and the base must be added to the solvent. A fact which merits consideration in any theory of acids and bases is that the speeds of large numbers of chemical reactions are greatly accelerated by acids and bases. This phenomenon is called acid-base catalysis, and a major reason for its wide prevalence is that most proton transfers are themselves exceedingly fast. Hence, reversible acid-base equilibria can usually be established very rapidly, and the resulting conjugate acids (or bases) then frequently offer favorable paths for the overall chemical reaction. The mechanisms of many of these catalyzed reactions are known. Some of them are specifically catalyzed by solvated protons (hydrogen ions); others, by hydroxide ions. Still others are catalyzed by acids or bases in the most general sense of the Br¨ onsted definitions. The existence of this general acid and base catalysis constituted an important item in the wide acceptance of the Br¨ onsted definitions. Lewis theory. Studies of catalysis have, however, played a large role in the acceptance of a set of quite different definitions of acids and bases, those due to G. N. Lewis. These definitions were originally proposed at about the same time as those of Br¨ onsted, but it was not until Lewis restated them in 1938 that they began to gain wide consideration. The Lewis definitions are: an acid is a substance which can accept an electron pair from a base; a base is a substance which can donate an electron pair. (These definitions are very similar to the terms popularized around 1927 by N. V. Sidgwick and others: electron donors, which are essentially Lewis bases, and electron acceptors, which are Lewis acids.) Bases under the Lewis definition are very similar to those defined by Br¨ onsted, but the Lewis definition for acids is very much broader. For example, virtually every cation is an acid, as are such species as AlCl3, BF3, and SO3. An acid-base reaction now typically becomes a combination of an acid with a base, rather than a proton transfer. Even so, many of the types of reactions which are characteristic of proton acids also will occur between Lewis acids and bases, for example, neutralization and color change of indicators as well as acid-base catalysis. Furthermore, these new definitions have been useful in suggesting new interrelations and in predicting new reactions, particularly for solid systems and for systems in nonaqueous solvents. For several reasons, these definitions have not been universally accepted. One reason is that the terms electron donor and electron acceptor had been widely accepted and appear to serve similar predicting and classifying purposes. A more important reason is unwillingness to surrender certain advantages in precision and definiteness inherent in the narrower Br¨ onsted definitions. It is a drawback of the Lewis definitions that the relative strengths of Lewis acids vary widely with choice of base and vice versa. For example, with the Br¨ onsted definitions, hydroxide ion is always a stronger base than ammonia;
with the Lewis definitions, hydroxide ion is a much weaker base than ammonia when reacting with silver ion but is stronger than ammonia when reacting with hydrogen ion. Another feature of the Lewis definitions is that some substances which have long been obvious examples of acid, for example, HCl and H2SO4, do not naturally fit the Lewis definition since they cannot plausibly accept electron pairs. Certain other substances, for example, carbon dioxide, are included by calling them secondary acids. These substances, too, tend to have electronic structures in which the ability to accept electron pairs is not obvious, but the more important distinction between them and primary acids is that their rates of neutralization by bases are measurably slow. However, in spite of these difficulties, the use of the Lewis definitions is increasing. Since there does not appear to be any simultaneous tendency to abandon the Br¨ onsted definitions, chemistry seems to be entering a period when the term acid needs a qualifying adjective for clarity, for example, Lewis acid or proton acid. Hard and soft acids and bases. As pointed out above, one of the drawbacks of such a broad definition as the Lewis one is that it is difficult to systematize the behavior of acids and bases toward each other. Attempts have been made to classify Lewis acids and bases into categories with respect to their mutual behavior. R. G. Pearson in 1963 proposed a simple and lucid classification scheme, based in part on earlier methods, that appears to be promising in its application to a wide variety of Lewis acid-base behavior. Lewis bases (electron donors) are classified as soft if they have high polarizability, low electronegativity, are easily oxidized, or possess low-lying empty orbitals. They are classified as hard if they have the opposite tendencies. Some bases, spanning the range of hard to soft and listed in order of increasing softness, are H2O, OH−, OCH3−, F−, NH3, C5H5N, NO2−, NH2OH, N2H4, C6H5SH, Br−, I−, SCN−, SO32−, S2O32−, and (C6H5)3P. Acids are divided more or less distinctly into two categories, hard and soft, with a few intermediate cases. Hard acids are of low polarizability, small size, and high positive oxidation state, and do not have easily excitable outer electrons. Soft acids have several of the properties of high polarizability, low or zero positive charge, large size, and easily excited outer electrons, particularly d electrons in metals. Some hard acids are H+, Li+, Na+, K+, Be2+, Mg2+, Ca2+, Sr2+, Mn2+, Al3+, Sc3+, Cr3+, Co3+, Fe3+, As3+, Ce3+, Si4+, Ti4+, Zr4+, Pu4+, BeMe2 (Me is the methyl group), BF3, BCl3, B(OR)3, Al(CH3)3, AlH3, SO3, and CO2. Examples of soft acids are Cu+, Ag+, Au+, Tl+, Hg+, Cs+, Pd2+, Cd2+, Pt2+, Hg2+, CH3Hg+, Tl3+, BH3, CO(CN)52−, I2, Br2, ICN, chloranil, quinones, tetracyanoethylene, O, Cl, Br, I, N, metal atoms, and bulk metals. Intermediate acids are Fe2+, Co2+, Ni2+, Cu2+, Pb2+, Sn2+, B(CH3)3, SO2, NO+, and R3C+. The rule for correlating acid-base behavior is then stated as follows: hard acids prefer to associate with hard bases and soft acids with soft bases. Application, for example, to the problem of OH− and NH3, mentioned earlier (and recognizing that OH− is hard
Acid and base compared with NH3), leads to reaction (14), which is unfavorable compared with reaction (15). OH− (hard) + Ag+ (soft) → AgOH
(14)
NH3 (soft) + Ag+ (soft) → Ag(NH3 )+
(15)
However, reaction (16) is unfavorable compared with reaction (17), which is in agreement with experiment. NH3 (soft) + H+ (hard) → NH4 +
(16)
OH− (hard) + H+ (hard) → H2 O
(17)
The rule is successful in correlating general acidbase behavior of a very wide variety of chemical systems, including metal-ligand interaction, chargetransfer complex formation, hydrogen bond formation, complex ion formation, carbonium ion formation, covalent bond formation, and ionic bond formation. It is to be emphasized that the hard-and-soft acidand-base concept is a means of classification and correlation and is not a theoretical explanation for acidbase behavior. The reasons why hard acids prefer hard bases and soft prefer soft are complex and varied. The already well-developed concepts of ionic and covalent bonds appear to be helpful, however. Hard acids and hard bases with small sizes and high charge would be held together with stable ionic bonds. Conversely, the conditions for soft acids and soft bases would be favorable for good covalent bonding. Existing theories of π-bonding also fit into the scheme. Usanovich theory. Another comprehensive theory of acids and bases was proposed by M. Usanovich in 1939 and is sometimes known as the positivenegative theory. Acids are defined as substances which form salts with bases, give up cations, and add themselves to anions and to free electrons. Bases are similarly defined as substances which give up anions or electrons and add themselves to cations. Two examples of acid-base reactions under this scheme are reactions (18) and (19). In the first, SO3 is an acid Na2 O + SO3 → 2Na+ + SO4 2−
(18)
2Na + Cl2 → 2Na+ + 2Cl−
(19)
because it takes up an anion, O2−, to form SO42−. In the second example, Cl2 is an acid because it takes up electrons to form Cl−. Using conventional terminology, this second reaction is an obvious example of oxidation-reduction. The fact that oxidationreduction can also be included in the Usanovich scheme is an illustration of the extensiveness of these definitions. So far, this theory has had little acceptance, quite possibly because the definitions are too broad to be very useful. Generation of chemical species. The acidic or basic character of a solvent can be used to stabilize interesting chemical species, which would otherwise be difficult to obtain. For example, carbonium ions have been thought to be important intermediates in
many organic reactions, but because of their fleeting existence as intermediates, their properties have been difficult to study. Most carbonium ions are very strong bases and would, for example, react, as shown by reaction (20). Accordingly, the equilibrium would R3 C+ + H2 O → R3 COH + H+
(20)
lie far to the right. However, use of a very strongly acidic solvent reverses the reaction, and measurable amounts of carbonium ion are then found. Concentrated sulfuric acid has found use in this connection. The very high acidity of SbF5 by itself, as a Lewis acid, and in mixtures with other Lewis acids, such as SO2, or protonic acids, such as HF and FSO3H, makes possible the study of many otherwise unstable carbonium ions. See SUPERACID. Acidity functions. A very different approach to the definition of acids, or perhaps better, to the definition of acidity, is to base the definition on a particular method of measurement. (As one example, it is probably true that the most nearly exact definition of pH is in terms of the electromotive force of a particular kind of galvanic cell.) It is possible to define various acidity functions in this way, and several have been proposed. One of the earliest and also one of the most successful is the H0 acidity function of L. P. Hammett. This defines an acidity in terms of the observed indicator ratio for a particular class of indicators, those which are uncharged in the basic form B. Suppose there is available a set of such indicators, and suppose further that the values of the dissociation constants of the acid forms BH+ are known. Then the h0 acidity of a solution is defined as Eq. (21), h0 = KBH+
[BH+ ] [B]
(21)
where KBH+ is the dissociation constant for the particular indicator employed, and where [BH+]/[B] is the experimentally observed ratio of concentrations of the conjugate acid and conjugate base forms of the indicator. To have a logarithmic scale (analogous to pH), the further definition is expressed in Eq. (22). H0 = − log h0
(22)
The virtues of this scale are that measurements are relatively simple and can be made for concentrated solutions and for solutions in nonaqueous or mixed solvents, situations where the pH scale offers difficulties. A further point is that in dilute aqueous solutions this new acidity becomes identical to pH. Although it has been found that this measure of acidity is fairly consistent within a class of indicators used, different classes can give somewhat different measures of acidity. Hence, caution must be used in interpretation of acidity measured by this technique. For a discussion of measurement of acidity see BASE (CHEMISTRY); BUFFERS (CHEMISTRY); HYDROGEN ION; IONIC EQUILIBRIUM; OXIDATION-REDUCTION; SOLUTION; SOLVENT. Franklin A. Long; Richard H. Boyd Bibliography. J. F. Coetzee and C. D. Ritchie, SoluteSolvent Interactions, vol. 1, 1969, vol. 2, 1976; C. W.
63
64
Acid-base indicator Hand, Acid-Base Chemistry, 1986; H. F. Holtzclaw, Jr., W. R. Robinson, and J. D. Odom, General Chemistry, 9th ed., 1991.
Acid-base indicator A substance that reveals, through characteristic color changes, the degree of acidity or basicity of solutions. Indicators are weak organic acids or bases which exist in more than one structural form (tautomers) of which at least one form is colored. Intense color is desirable so that very little indicator is needed; the indicator itself will thus not affect the acidity of the solution. The equilibrium reaction of an indicator may be regarded typically by giving it the formula HIn. It dissociates into H+ and In− ions and is in equilibrium with a tautomer InH which is either a nonelectrolyte or at most ionizes very slightly. In the overall equilibrium shown as reaction (1), the simplifying assumplnH
H+ + ln−
Hln ionization
(1)
dissociation
(1)
(2)
tion that the indicator exists only in forms (1) and (2) leads to no difficulty. The addition of acid will completely convert the indicator to form (1), which is therefore called the acidic form of the indicator although it is functioning as a base. A hydroxide base converts the indicator to form (2) with the formation of water; this is called the alkaline form. For the equilibrium between (1) and (2) the equilibrium
constant is given by Eq. (2). In a manner similar to + − H In KIn = (2) [InH] the pH designation of acidity, that is, pH = −log[H+], the KIn is converted to pKIn with the result shown in Eq. (3). It is seen that the pK of an indicator has pKIn = pH − log
[In− ] [InH]
(3)
a numerical value approximately equal to that of a specific pH level. Use of indicators. Acid-base indicators are commonly employed to mark the end point of an acidbase titration or to measure the existing pH of a solution. For titration the indicator must be so chosen that its pK is approximately equal to the pH of the system at its equivalence point. For pH measurement, the indicator is added to the solution and also to several buffers. The pH of the solution is equal to the pH of that buffer which gives a color match. Care must be used to compare colors only within the indicator range. A color comparator may also be used, employing standard color filters instead of buffer solutions. Indicator range. This is the pH interval of color change of the indicator. In this range there is competition between indicator and added base for the available protons; the color change, for example, yellow to red, is gradual rather than instantaneous. Observers may, therefore, differ in selecting the precise point of change. If one assumes arbitrarily that the indicator is in one color form when at least 90% of it is in that form, there will be uncertain color
Common acid-base indicators Common name
pH range
Color change (acid to base)
Methyl violet
0 2, 5 6
Yellow to blue violet to violet
Metacresol purple
1.2 2.8, 7.3 9.0
Red to yellow to purple
Thymol blue
1.2 2.8, 8.0 9.6
Red to yellow to blue
Tropeoline 00 (Orange IV) Bromphenol blue
1.4 3.0
Red to yellow
3.0 4.6
Methyl orange
pK
Chemical name
Structure
Solution
Base
0.25% in water
1.5
Pentamethylbenzyl-pararosaniline hydrochloride m-Cresolsulfonphthalein
Acid
1.7
Thymolsulfonphthalein
Acid Base
Yellow to blue
4.1
Sodium p-diphenyl-aminoazobenzenesulfonate Tetrabromophenol-sulfonphthalein
0.73% in N/50 NaOH, dilute to 0.04% 0.93% in N/50 NaOH, dilute to 0.04% 0.1% in water
2.8 4.0
Orange to yellow
3.4
Base
Bromcresol green Methyl red
3.8 5.4 4.2 6.3
Yellow to blue Red to yellow
4.9 5.0
Acid Base
Chlorphenol red
5.0 6.8
Yellow to red
6.2
Sodium p-dimethyl-aminoazobenzenesulfonate Tetrabromo-m-cresolsulfonphthalein Dimethylaminoazobenzene-o-carboxylic acid Dichlorophenolsulfonphthalein
Bromcresol purple
5.2 6.8
Yellow to purple
6.4
Dibromo-o-cresolsulfonphthalein
Acid
Bromthymol blue
6.0 7.6
Yellow to blue
7.3
Dibromothymolsulfonphthalein
Acid
Phenol red
6.8 8.4
Yellow to red
8.0
Phenolsulfonphthalein
Acid
Cresol red
2.0 3.0, 7.2 8.8
Orange to amber to red
8.3
o-Cresolsulfonphthalein
Acid
Orthocresolphthalein Phenolphthalein Thymolphthalein Alizarin yellow GG Malachite green
8.2 9.8 8.4 10.0 10.0 11.0 10.0 12.0 11.4 13.0
Colorless to red Colorless to pink Colorless to red Yellow to lilac Green to colorless
— — — Sodium p-nitrobenzeneazosalicylate p,p -Benzylidene-bis-N,N-dimethylaniline
Acid Acid Acid Acid Base
9.7 9.9
Acid
Acid
1.39% in N/50 NaOH, dilute to 0.04% 0.1% in water 0.1% in 20% alcohol 0.57% in N/50 NaOH, dilute to 0.04% 0.85% in N/50 NaOH, dilute to 0.04% 1.08% in N/50 NaOH dilute to 0.04% 1.25% in N/50 NaOH, dilute to 0.04% 0.71% in N/50 NaOH, dilute to 0.04% 0.76% in N/50 NaOH, dilute to 0.04% 0.04% in alcohol 1% in 50% alcohol 0.1% in alcohol 0.1% in warm water 0.1% in water
Acid anhydride in the range of 90–10% InH (that is, 10–90% In−). When these arbitrary limits are substituted into the pK equation, the interval between definite colors is shown by Eqs. (4). Thus the pH uncertainty of the pH = pK + log
10 90
to
pH = pK + log
90 10
(4)
indicator is from pK + 1 to pK − 1 (approximately), and pK ± 1 is called the range of the indicator. The experimentally observed ranges may differ from this prediction somewhat because of variations in color intensity. Examples. The table lists many of the common indicators, their chemical names, pK values, ranges of pH, and directions for making solutions. Many of the weak-acid indicators are first dissolved in N/50 NaOH to the concentration shown, then further diluted with water. The weak-base indicators show some temperature variation of pH range, following approximately the temperature change of the ionization constant of water. Weak-acid indicators are more stable toward temperature change. The colors are formed by the usual chromophoric groups, for example, quinoid and azo-. One of the most common indicators is phenolphthalein, obtained by condensing phthalic anhydride with phenol. The acid form is colorless. It is converted by OH− ion to the red quinoid form, as shown by reaction (5). OH
O
O CH3
C
CH3
OH
C
O O
C
CH3
(2)
(1)
Anhydrides of straight-chain acids containing from 2 to 12 carbon atoms are liquids with boiling points higher than those of the parent acids. They are relatively insoluble in cold water and are soluble in alcohol, ether, and other common organic solvents. The lower members are pungent, corrosive, and weakly lacrimatory. Anhydrides from acids with more than 12 carbon atoms and cyclic anhydrides from dicarboxylic acids are crystalline solids. Preparation. Because the direct intermolecular removal of water from organic acids is not practicable, anhydrides must be prepared by means of indirect processes. A general method involves interaction of an acid salt with an acid chloride, reaction (1). O CH3
O O– NA+ + CH 3
C
C
CI NaCI + (CH 3CO)2O
(1)
Acetic anhydride, the most important aliphatic anhydride, is manufactured by air oxidation of acetaldehyde, using as catalysts the acetates of copper and cobalt, shown in reaction (2); peracetic acid appar-
OH
O catalyst
H + O2
CH3 C
O
C NaOH
O
CH3 COOH + CH 3CHO
(CH3CO)2O + H 2O
(2)
C O OH
C C
ently is an intermediate. The anhydride is separated from the by-product water by vacuum distillation. Another important process utilizes the thermal decomposition of ethylidene acetate (made from acetylene and acetic acid), reaction (3).
O
(5) ONa
mercury salt
2CH3COOH + HC OOCCH3
O
Other indiators such as methyl orange and methyl red are sodium salts of azobenzene derivatives containing sulfonic and carboxylic groups respectively. See ACID AND BASE; HYDROGEN ION; TITRATION. Allen L. Hanson
CH3
C
CH catalyst
heat
OOCCH3
ZnCI2
One of an important class of reactive organic compounds derived from acids via formal intermolecular dehydration; thus, acetic acid (1), on loss of water, forms acetic anhydride (2).
(3)
Acetic anhydride has been made by the reaction of acetic acid with ketene, reaction (4). CH3COOH + CH 2
Acid anhydride
CH3CHO + (CH 3CO)2O
H
C
O
(CH3CO)2O
(4)
Mixed anhydrides composed of two different radicals are unstable, and disproportionate to give the two simple anhydrides. Direct use is made of this in the preparation of high-molecular-weight anhydrides, as seen in reaction (5). The two simple anhydrides are easily separable by distillation in a vacuum.
65
66
Acid halide
(5)
Cyclic anhydrides are obtained by warming succinic or glutaric acids, either alone, with acetic anhydride, or with acetyl chloride. Under these conditions, adipic acid first forms linear, polymeric anhydride mixtures, from which the monomer is obtained by slow, high-vacuum distillation. Cyclic anhydrides are also formed by simple heat treatment of cis-unsaturated dicarboxylic acids, for example, maleic and glutaconic acids; and of aromatic 1,2dicarboxylic acids, for example, phthalic acid. Commercially, however, both phthalic (3) and maleic (4) O
O HC
C
O HC
C
O (3)
Acid halide
C
O C
fluorescein. Phthalic anhydride is used in manufacturing glyptal resins (from the anhydride and glycerol) and in manufacturing anthraquinone. Heating phthalic anhydride with ammonia gives phthalimide, used in Gabriel’s synthesis of primary amines, amino acids, and anthranilic acid (o-aminobenzoic acid). With alkaline hydrogen peroxide, phthalic anhydride yields monoperoxyphthalic acid, used along with benzoyl peroxide as polymerization catalysts, and as bleaching agents for oils, fats, and other edibles. Anhydrides react with water to form the parent acid, with alcohols to give esters, and with ammonia to yield amides; and with primary or secondary amines, they furnish N-substituted and N,Ndisubstituted amides, respectively. See ACID HALIDE; ACYLATION; CARBOXYLIC ACID; DIELS-ALDER REACTION; ESTER; FRIEDEL-CRAFTS REACTION. Paul E. Fanta Bibliography. F. Bettelheim and J. March, Introduction to Organic and Biochemistry, 1997; R. J. Fessenden and J. S. Fessenden, Organic Chemistry, 6th ed., 1998.
One of a large group of organic substances possessing the halocarbonyl group O
O (4)
anhydrides are primary products of manufacture, being formed by vapor-phase, catalytic (vanadium pentoxide), air oxidation of naphthalene and benzene, respectively; at the reaction temperature, the anhydrides form directly. Uses. Anhydrides are used in the preparation of esters. Ethyl acetate and butyl acetate (from butyl alcohol and acetic anhydride) are excellent solvents for cellulose nitrate lacquers. Acetates of highmolecular-weight alcohols are used as plasticizers for plastics and resins. Cellulose and acetic anhydride give cellulose acetate, used in acetate rayon and photographic film. The reaction of anhydrides with sodium peroxide forms peroxides (acetyl peroxide is violently explosive), used as catalysts for polymerization reactions and for addition of alkyl halides to alkenes. In Friedel-Crafts reactions, anhydrides react with aromatic compounds, forming ketones such as acetophenone. Maleic anhydride reacts with many dienes to give hydroaromatics of various complexities (Diels-Alder reaction). Maleic anhydride is used commercially in the manufacture of alkyd resins from polyhydric alcohols. Soil conditioners are produced by basic hydrolysis of the copolymer of maleic anhydride with vinyl acetate. Phthalic anhydride and alcohols form esters (phthalates) used as plasticizers for plastics and resins. Condensed with phenols and sulfuric acid, phthalic anhydride yields phthaleins, such as phenolphthalein; with m-dihydroxybenzenes under the same conditions, xanthene dyes form, for example,
R
C
X
in which X stands for fluorine, chlorine, bromine, or iodine. The terms acyl and aroyl halides refer to aliphatic or aromatic derivatives, respectively. The great inherent reactivity of acid halides precludes their free existence in nature; all are made by synthetic processes. In general, acid halides have low melting and boiling points and show little tendency toward molecular association. With the exception of the formyl halides (which do not exist), the lower members are pungent, corrosive, lacrimatory liquids that fume in moist air. The higher members are lowmelting solids. Preparation. Acid chlorides are prepared by replacement of carboxylic hydroxyl of organic acids by treatment with phosphorus trichloride, phosphorus pentachloride, or thionyl chloride [reactions (1)– (3)]. 3RCOOH + PCl3 → 3RCOCl + H3 PO3
(1)
RCOOH + PCl5 → RCOCl + POCl3 + HCl
(2)
RCOOH + SOCl2 → RCOCl + SO2 + HCl
(3)
Although acid bromides may be prepared by the above methods (especially by use of PBr3), acid iodides are best prepared from the acid chloride treatment with either CaI2 or HI, and acid fluorides from the acid chloride by interaction with HF or antimony fluoride. Reactions and uses. The reactivity of acid halides centers upon the halocarbonyl group, resulting in substitution ofthe halogen by appropriate
Acid rain structures. Thus, as shown by reactions (4), with
HOH → HCl + RCOOH C2 H5 OH → HCl + RCOOC2 H5 RCOCl + 2NH3 → NH4 Cl + RCONH2 2HNR2 → NH2 R2 Cl + RCONR2
(4)
substances containing active hydrogen atoms (for example, water, primary and secondary alcohols, ammonia, and primary and secondary amines), hydrogen chloride is formed together with acids, esters, amides, and N-substituted amides, respectively. The industrially prepared acetyl and benzoyl chlorides are much used in reactions of the above type, particularly in the acetylating or benzoylating of amines and amino acids, and alcohols, especially polyalcohols such as glycerol, the sugars, and cellulose, to form amides, esters, and polyesters, respectively. In reactions (4) the by-product hydrogen chloride must be neutralized. With aliphatic acid halides, this must be done by using an excess of the amine (as shown); aromatic acid chlorides, being insoluble in water, can be used in the presence of aqueous sodium hydroxide (Schotten-Baumann reaction). An alternative technique uses the tertiary amine, pyridine, both as solvent and as neutralizing agent. Interaction of acid chlorides with sodium salts of organic acids furnishes a general method for the preparation of acid anhydrides, as shown by reaction (5). Analogously, aromatic acyl peroxides, used RCOCl + RCOONa → NaCl + (RCO)2 O
(5)
as bleaching agents for flour, fats, and oils, and as polymerization catalysts, may be prepared from sodium peroxide and the acid halide [reaction (6)]. 2C6 H5 COCl + Na2 O2 → 2NaCl + C6 H5 COOOOCC6 H5 Benzoyl peroxide
(6)
See PEROXIDE. Acid halides are used to prepare ketones, either by reaction with alkyl or aryl cadmium reagents [reaction (7)] or by the aluminum chloride catalyzed Friedel-Crafts reaction [reaction (8)]. RCOCl + RCH2 CdCl → RCOCH2 R + CdCl2 AlCl3
RCOCl + C6 H6 −−−−−−→ RCOC6 H5 + HCl
(7) (8)
Reduction of acid halides is easily effected. In the Rosenmund method, used mainly for aromatic acid halides, hydrogen and a poisoned palladium catalyst are employed [reaction (9)]. The product here Pd
RCOCl + H2 −−−−−−→ HCl + RCHO poisoned
(9)
is the aldehyde. For reduction to the alcohol stage, the vigorous reagent lithium aluminum hydride reduces both aliphatic and aromatic acid halides [reaction (10)]. 4HCl
4RCOCl + LiAlH4 → LiAlCl4 + (RCH2 O)4 LiAl −−−−−−→ 4RCH2 OH + LiAlCl4 (10)
Direct substitution of halogen (chlorine fastest, bromine more slowly) into acid aliphatic halides is relatively easy, compared with substitution in the parent acid, and takes place on the carbon next to the carbonyl group. The product is an α-halo acid halide, RCHXCOX. These compounds interact with carboxylic acids, via an equilibrium, to form an α-halo acid and an acid halide [reaction (11)]. RCHXCOX + RCH2 COOH
RCHXCOOH + RCH2 COX
(11)
Thus, if a small amount of acyl halide is added to a large amount of carboxylic acid, the latter can be chlorinated or brominated to completion by treatment with either chlorine or bromine (Hell-VolhardZelinski reaction). See ACYLATION; CARBOXYLIC ACID; FRIEDEL-CRAFTS REACTION. Paul E. Fanta Bibliography. N. L. Allinger et al., Organic Chemistry, 2d ed., 1976; W. H. Brown, Introduction to Organic and Biochemistry, 4th ed., 1987.
Acid rain Precipitation that incorporates anthropogenic acids and acidic materials. The deposition of acidic materials on the Earth’s surface occurs in both wet and dry forms as rain, snow, fog, dry particles, and gases. Although 30% or more of the total deposition may be dry, very little information that is specific to this dry form is available. In contrast, there is a large and expanding body of information related to the wet form: acid rain or acid precipitation. Acid precipitation, strictly defined, contains a greater concentration of hydrogen (H+) than of hydroxyl (OH−) ions, resulting in a solution pH less than 7. Under this definition, nearly all precipitation is acidic. The phenomenon of acid deposition, however, is generally regarded as resulting from human activity. See PRECIPITATION (METEOROLOGY). Sources. Theoretically, the natural acidity of precipitation corresponds to a pH of 5.6, which represents the pH of pure water in equilibrium with atmospheric concentrations of carbon dioxide. Atmospheric moisture, however, is not pure, and its interaction with ammonia, oxides of nitrogen and sulfur, and windblown dust results in a pH between 4.9 and 6.5 for most “natural” precipitation. The distribution and magnitude of precipitation pH in the United States (Fig. 1) suggest the impact of anthropogenic rather than natural causes. The areas of highest precipitation acidity (lowest pH) correspond to areas within and downwind of heavy industrialization and urbanization where emissions of sulfur and nitrogen oxides are high. It is with these emissions that the most acidic precipitation is thought to originate. See ACID RAIN. Atmospheric processes. The transport of acidic substances and their precursors, chemical reactions, and deposition are controlled by atmospheric processes. In general, it is convenient to distinguish between physical and chemical processes, but it must
67
68
Acid rain
4.2 5.0
4.5
5.5 4.2 4.5
Fig. 1. Distribution of rainfall pH in the eastern United States.
be realized that both types may be operating simultaneously in complicated and interdependent ways. The physical processes of transport by atmospheric winds and the formation of clouds and precipitation strongly influence the patterns and rates of acidic deposition, while chemical reactions govern the forms of the compounds deposited. In midlatitude continental regions, such as eastern North America, most precipitation arises from cyclonic storms. As shown by the schematic represen-
acid rain upper-level outflow
L
warm front
low inf level low
ac id
ra in
warm sector
N
photochemical reactions
cold front
Fig. 2. Schematic view of a typical cyclonic storm in eastern North America with a lowpressure region (L).
tation of a typical storm (Fig. 2), rain tends to form along the surface cold and warm fronts that define the so-called warm sector of the storm. This characteristic structure, arising from prevailing north-south temperature gradients, simultaneously sets the stage for chemical transformations of pollutants and for incorporation of the compounds into precipitation. The motion of the cold front toward the southeast forces the moist warm-sector air to stream along the cold front toward the low-pressure region (L), as shown by the broad arrow in Fig. 2. At the same time, the air is gradually lifted out of the surface layer and into the colder, upper parts of the storm, which allows the supply of water vapor to condense out, forming the cloud and precipitation. See CYCLONE; FRONT; STORM. There are a number of chemical pathways by which the primary pollutants, sulfur dioxide (SO2) from industry, nitric oxide (NO) from both industry and automobiles, and reactive hydrocarbons mostly from trees, are transformed into acid-producing compounds. Some of these pathways exist solely in the gas phase, while others involve the aqueous phase afforded by the cloud and precipitation. As a general rule, the volatile primary pollutants must first be oxidized to more stable compounds before they are efficiently removed from the atmosphere. Ironically, the most effective oxidizing agents, hydrogen peroxide (H2O2) and ozone (O3), arise from photochemical reactions involving the primary pollutants themselves. See AIR POLLUTION. All of the ingredients needed to form the strong mineral acids of sulfur and nitrogen [sulfuric acid (H2SO4) and nitric acid (HNO3)], which constitute most of the acids found in rain, exist in the warm sector. Especially in summertime, stagnant air conditions trap the pollutants under clear skies for several
Acid rain days, permitting photochemical reactions to initiate other gas-phase reactions. Nitrogen oxides (NOx) in combination with ultraviolet light from the Sun, reactive hydrocarbons, atmospheric oxygen, and water vapor simultaneously give rise to HNO3 vapor, odd hydrogen radicals (particularly OH), and the strong oxidants H2O2 and O3. Meanwhile, slow oxidation of sulfur dioxide (SO2), initiated by reaction with the OH radical, leads to the gradual buildup of sulfate (SO−24) aerosol, and hazy skies result from the highly condensable nature of H2SO4. While some dry deposition does occur under these conditions, the warm sector of midlatitude cyclones is a region conducive to the gradual accumulation of primary pollutants, oxidants, and acid-forming compounds in the lower atmosphere. See SMOG. As the cold front approaches any given location in the warm sector, the airborne acids and their precursors are drawn into the circulations of the cyclone, particularly into the stream of moist air associated with the fronts. By this large-scale meteorological process, the pollutants become integrally associated with the frontal cloud systems, and attach to individual cloud particles by a variety of microphysical processes. Since most atmospheric sulfate particles are very soluble, they act as good centers (nuclei) for cloud drop formation. Thus, the process of cloud formation is itself a mechanism for scavenging particulate pollutants from the air. As the cloud drops grow by condensation, soluble gases will be absorbed by the cloud water. It is during this phase of the overall process that additional SO2 will be oxidized to H2SO4 by the previously formed strong oxidants, particularly the H2O2. The actual removal of pollutants from the atmosphere by wet deposition requires the formation of precipitation within the clouds. Without cloud elements greater than about 100 micrometers in diameter, the pollutant mass remains in the air, largely in association with the relatively small cloud drops, which have negligible rates of descent. Since most precipitation in midlatitude storms is initiated by the formation of ice particles in the cold, upper reaches of the clouds, pollutant fractionation between water phases inhibits the transfer of cloud water acidity to large particles, which do precipitate readily. Partly because of such microphysical phenomena and partly because of nonuniform distributions of pollutants within the clouds, the acidity of precipitation tends to be substantially less than that of the cloud water that remains aloft. The acidity of the precipitation appears to be acquired largely through collisions of the melted ice particles with the relatively concentrated cloud drops in the lower portions of the storm clouds. The rate of wet deposition is thus governed by a complicated set of meteorological and microscale processes working in harmony to effect a general cleansing of the atmosphere. See ATMOSPHERIC CHEMISTRY; CLOUD PHYSICS. Terrestrial and aquatic effects. The effect of acid deposition on a particular ecosystem depends largely on its acid sensitivity, its acid neutralization capabil-
ity, the concentration and composition of acid reaction products, and the amount of acid added to the system. As an example, the major factors influencing the impact of acidic deposition on lakes and streams are (1) the amount of acid deposited; (2) the pathway and travel time from the point of deposition to the lake or stream; (3) the buffering characteristics of the soil through which the acidic solution moves; (4) the nature and amount of acid reaction products in soil drainage and from sediments; and (5) the buffering capacity of the lake or stream. In many ecosystems, except for foliar effects, the impacts of acid precipitation on aquatic and terrestrial ecosystems overlap. This is because soils are the key intermediate. They provide the root environment for terrestrial vegetation, and also control the water quality of runoff and soil drainage, which supplies most of the water to the aquatic system. A number of acid-consuming reactions occur in soil, which lessen the impact of acid additions on both the soil and soil drainage waters. In soils, indigenous or agriculturally amended carbonates (CO32−) react with acid precipitation to raise the pH of soil drainage waters, while maintaining soil pH. Also, soils exhibit a cation exchange capacity that can serve to neutralize acid precipitation. At neutral soil pH (6.0–8.0), most cations on the exchange are calcium (Ca) and magnesium (Mg). When acids are added, H+ from solution exchanges with the adsorbed Ca2+ and Mg2+. Although this reduces the acidity of soil drainage waters, the soil acidity is increased. Eventually, when the neutralizing carbonates and exchangeable calcium and magnesium supplies are exhausted, soil minerals react with the acid. At a soil pH of about 5.2 or less, substantial aluminum (Al) can be solubilized from the dissolution of clay minerals and coatings on soil particles. See SOIL CHEMISTRY. Acid deposition directly into lakes and streams causes chemical reactions analogous to those in the soil system. However, instead of soil materials, the carbonate-bicarbonate (HCO−3) system buffers the solution. As with soils, waters of low pH and buffering capacities, in this case bicarbonate concentrations, are most sensitive. At solution pH of 4.5, bicarbonate is essentially depleted, and subsequent acid additions that reduce pH also mobilize metals from suspended solids and the lake or stream bed. Generally, lakes and streams that drain acidsusceptible soils are most susceptible to direct acid deposition as well. Shallow soils with low pH, low cation-exchange capacity, and high permeability not only are least able to neutralize acid flows but also are poor sources of the bicarbonate waters needed to buffer the aquatic system. See BUFFERS (CHEMISTRY). While the study of acid precipitation effects on terrestrial ecosystems is relatively new, soil acidification is a naturally occurring process in humid climates and has long been the subject of research, whose findings suggest acid precipitation effects. The generally accepted impact of soil acidification on the productivity of terrestrial plants is summarized as follows. As soil becomes more acidic, the basic cations (Ca, Mg) on the soil exchange are replaced by
69
70
Acid rain H2SO4, HNO3 (sulfuric, nitric acids) disassociation
H
K+
+
+
mobile anions move in the soil solution with displaced and solubilized nutrient cations
Ca2+
Ca2+ soil with exchangeable nutrient cations
H+
2− − SO4 , NO3
H+ Mg2+
Mg2+
Ca2+
Ca2+
Ca2+, Mg2+, K+, etc. +SO42−, NO3−, Al3+, Ni, Zn, Cu, Fe, Mn
displacement of nutrient cations results in leached soil of increased acidity H+
Al3+ H+
nutrient ions lost to drainage waters
Al 3+
soil leached of nutrient cations
Al3+
H+ +
Al3 H+
Al3+
Ca2+
aquatic ecosystems, lakes, streams
Fig. 3. Soil acidification and loss of soil nutrients following acid inputs.
hydrogen ions or solubilized metals. The basic cations, now in solution, can be leached through the soil (Fig. 3). As time progresses, the soil becomes less fertile and more acidic. Resultant decreases in soil pH cause reduced, less-active populations of soil microorganisms, which in turn slow decomposition of plant residues and cycling of essential plant nutrients. The plant availability of phosphorus, now composing mostly aluminum and iron (Fe) phosphates, decreases, while the plant availability of trace metals [aluminum, copper (Cu), iron, zinc (Zn), boron (B), manganese (Mn)] increases, sometimes to phytotoxic levels. Acid precipitation may injure trees directly or indirectly through the soil. Foliar effects have been studied extensively, and it is generally accepted that visible damage occurs only after prolonged exposure to precipitation of pH 3 or less (for example, acid fog or clouds). Measurable effects on forest ecosystems will then more likely result indirectly through soil processes than directly through exposure of the forest canopy. The characteristics that make a soil sensitive to acidification are extremely variable from soil to soil. Accordingly, forest decline will not occur regionally across a deposition gradient, but it will occur in localized areas with sensitive soils. The occurrence of both acid-sensitive and acid-tolerant soils in areas receiving high levels of acid deposition and in areas
receiving low levels of acid deposition makes causeand-effect relationships between acid deposition and forest decline difficult to establish. For example, 10 years of monitoring indicated that hardwood forests in Ontario were on average healthy and stable. Still, tree condition was poorer in areas with both moderately high pollution deposition and acidsensitive soils. The clearest cause-and-effect relationship exists for red spruce and white birch exposed to acid clouds and fog. Areas of declining red spruce and white birch coincide with areas of acid clouds and fog, and the extent of decline is related to the duration of exposure. Many important declines in the condition of forest trees have been reported in Europe and North America during the period of increasing precipitation acidity. These cases include injury to white pine in the eastern United States, red spruce in the Appalachian Mountains of eastern North America, and many economically important species in central Europe. Since forest trees are continuously stressed by competition for light, water, and nutrients; by disease organisms; by extremes in climate; and by atmospheric pollutants, establishing acid deposition as the cause of these declines is made more difficult. Each of these sources of stress, singly or in combination, produces similar injury. However, a large body of information indicates that accelerated soil acidification resulting from acid deposition is an important predisposing stress that in combination with other stresses has resulted in increased decline and mortality of sensitive tree species and widespread reduction in tree growth. See FOREST ECOSYSTEM; TERRESTRIAL ECOSYSTEM. Aquatic biology effects. Aquatic ecosystems are sustained by intricate exchanges of energy and material among organisms at all levels. Acidic deposition impacts aquatic ecosystems by harming individual organisms and by disrupting flows of energy and materials through the ecosystem. The effect of acid deposition is commonly assessed by studying aquatic invertebrates and fish. Aquatic invertebrates live in the sediments of lakes and streams and are vitally important to the cycling of energy and material in aquatic ecosystems. These small organisms, rarely larger than 20 mm (1 in.), break down large particulate organic matter for further degradation by microorganisms, and they are an important food source for fish, aquatic birds, and predatory invertebrates. Individual aquatic organisms are harmed through direct physiological effects caused by low pH, high metal concentrations, and the interaction between pH and metals. Low pH causes lesions to form on gills and erodes gill tissue. Gill damage compromises respiration, excretion, and liver function, which in turn leaves organisms susceptible to further health problems. As pH decreases, the aluminum concentration of streams and lakes increases. Aluminum enhances the toxic effects of low pH. Together these phenomena result in severe acidosis and cause gills to secrete excessive amounts of mucus, further
Acid rain interfering with respiration by reducing oxygen diffusion across gills. Additionally, gill membranes become more permeable, reducing sodium uptake and increasing sodium loss. The most likely physiological cause of fish death at pH 4–5 is the inability to maintain a proper salt balance (ion regulation). Metals such as aluminum liberated from acidified soils may accumulate externally on the gill surfaces of invertebrates, and this accumulation also impairs respiration. Iron has been shown to accumulate on the gut membranes of invertebrates, inhibiting food absorption. It should be noted that the formation of metallo-organic complexes reduces the toxicity of metals in low-pH waters. Acid bogs and marshes support healthy populations of aquatic organisms because soluble metals and humic compounds form such complexes. Currently, there are concerns that acid deposition is causing the loss of fish species, through the physiological damage discussed above and by reproductive impairment. In order to find water with a suitable pH, fish may be forced to spawn in areas that are less than optimal for egg development, with a consequent reduction in the number of, or total loss of, offspring. Female fish need high serum calcium levels for normal ovary formation and egg laying. Female fish in acidified waters maintain lower than normal serum calcium levels, and consequently egg laying may be inhibited, again resulting in fewer or no offspring. Early life stages of fish (eggs and fry) are usually more sensitive than adults to low pH, and acute mortality of these life stages may be high during periods of rapid declines in pH following rainfall. The death of fish embryos is thought to result from the corrosion of skin cells, which interferes with respiration and osmoregulation. While fish die from acidification, their numbers and diversity
are more likely to decline from a failure to reproduce. The effects of acid deposition on individuals described above in turn elicit changes in the composition and abundance of communities of aquatic organisms. The degree of change depends on the severity of acidification, and the interaction of other factors, such as metal concentrations and the buffering capacity of the water. The pattern most characteristic of aquatic communities in acidified waters is a loss of species diversity, and an increase in the abundance of a few, acid-tolerant taxa. At moderately affected sites, acid-sensitive species are frequently replaced by more acid-tolerant forms, while the total number or mass of organisms may not change. For example, ephemeropterans (mayflies) are quite sensitive to reductions in pH and their numbers may decline, while plecopterans (stone flies) and dipterans (true flies) become relatively more dominant following acidification. Invertebrates in the feeding group that shred leaves also become relatively more dominant than those that graze on diatoms in acid-impacted streams. As pH decreases further and toxic metals increase, both the number of species and the number of organisms decrease sharply (Fig. 4). Mollusks and crustaceans are particularly sensitive to acidification, since both groups require high levels of calcium to form their shells. Calcium is scarce in the poorly buffered systems that are most affected by acidic deposition. Community-level effects may occur indirectly, as a result of changes in the food supply and in predatorprey relations. Reduction in the quality and amount of periphyton may decrease the number of herbivorous invertebrates, which may in turn reduce the number of organisms (predatory invertebrates and fish) that feed upon herbivorous invertebrates. The
8 12
3
inc 4
re a
sin
4
g
4
ab
un
4
da
nc
e
de
cre
ing
rs
ity
12
pH
Key: Ephemeroptera Plecoptera Diptera Others
di
ve
3
as
an
d
2
1 Trichoptera 6
Fig. 4. General pattern of aquatic invertebrate abundance and diversity with a reduction in pH. Total abundance is indicated by disk sizes. Numbers of species within each order are indicated.
71
72
Acipenseriformes disappearance of fish may result in profound changes in plant and invertebrate communities. Dominant fish species function as keystone predators, controlling the size distribution, diversity, and numbers of invertebrates. Their reduction alters the interaction within and among different levels of the food web and the stability of the ecosystem as a whole. Acid deposition affects the survival of individual organisms and the relationship of communities within aquatic ecosystems. Early life stages are generally more sensitive to low pH then adult organisms. Directs effects of pH and aluminum toxicity affect survival more than the indirect effects of food availability or predation. The rate and magnitude of acidification of aquatic systems depend upon the factors discussed above. Predictable changes in water chemistry, as the pH of surface waters decreases, include decreased alkalinity, decreased buffering capacity, and increased concentrations of aluminum, magnesium, and iron, among other elements. Such changes in water chemistry generally result in decreased diversity of aquatic species and reduced productivity. The impact of acid deposition on terrestrial and aquatic ecosystems is not uniform. While increases in acid deposition may stress some ecosystems and reduce their stability and productivity, others may be unaffected. The degree and nature of the impact depend on the acid input load, organismal susceptibility, and buffering capacity of the particular ecosystem. See BIOGEOCHEMISTRY; ECOSYSTEM. Ronald R. Schnabel; Dennis Lamb; Harry B. Pionke; Dennis Genito Bibliography. Acid Deposition: Atmospheric Processes in Eastern North America, National Academy Press, 1983; Acid Rain and Transported Air Pollutants: Implications for Public Policy, U.S. Congress, Office of Technology Assessment, 1984; R. J. Charlson and H. Rohde, Factors controlling the acidity of natural rainwater, Nature, 195:683–685, 1982; F. M. D’Itri (ed.), Acid Precipitation: Effects on Ecological Systems, 1982; D. L. Godbold and A. H¨ uttermann (eds.), Effects of Acid Rain on Forest Processes, Wiley-Liss, New York, 1994; R. M. Heard et al., Episodic acidification and changes in fish diversity in Pennsylvania headwater streams, Trans. Amer. Fisheries Soc., 126:977–984, 1997; Interim Assessment: The Causes and Effects of Acidic Deposition, National Acid Precipitation Assessment Program, Washington, D.C., 1987; W. G. Kimmel et al., Macroinvertebrate community structure and detritus processing rates in two southwestern Pennsylvania streams acidified by atmospheric deposition, Hydrobiologia, 124:97–102, 1985; W. P. Robarge and D. W. Johnson, The effects of acidic deposition on forested soils, Adv. Agron., 47:1–83, 1992; A. D. Rosemond et al., The effects of stream acidity on benthic invertebrate communities in the south-eastern United States, Freshwater Biol., 27:193–209, 1992; J. N. Wordman and E. B. Cowling, Airborne chemicals and forest health, Environ. Sci. Technol., 21:120– 126, 1987.
Acipenseriformes An order of actinopterygian fishes in the subclass Chondrostei, which probably originated from palaeonisciform ancestors during the Triassic or Jurassic periods. Acipenseriformes, the only extant order of about 10 orders of chondrosteans, are characterized by a largely cartilaginous skeleton, heterocercal tail, fins with more than one ray per pterygiophore, spiral valve intestine, myodome and preopercle reduced or absent, and gular plate absent. The cartilaginous skeleton of Acipenseriformes, long regarded as primitive, was taken as indicative of an alliance with the Chondrichthyes; however, cartilage is secondary and the true relationship is with bony fishes. See ACTINOPTERYGII; CHONDROSTEI; OSTEICHTHYES. Acipenseriformes includes two suborders: the fossil Chondrosteoidei, which became extinct in the Mesozoic, and the Acipenseroidei, with two families, Acipenseridae and Polyodontidae, both with little apparent change from the Cretaceous. Acipenseridae (sturgeons). Sturgeons are characterized by a robust body armed with five longitudinal rows of bony plates; a pronounced snout bearing sensitive barbels on the underside; an inferior mouth, which is protrusible and, in adults, toothless; spiracles present or absent; gill rakers fewer than 50; a large swim bladder; and pectoral fins, the leading edge of which consists of fused rays. The family includes four extant genera and about 25 species in two subfamilies. The subfamily Acipenserinae comprises three genera: Acipenser (see illustration) is distinguished by having spiracles and a subconical snout, and is known from North America, Europe, and Asia. In Pseudoscaphirhynchus, which is endemic to the Aral Sea basin, spiracles are absent, the snout is depressed, and the caudal peduncle is short, depressed, and incompletely armored. Scaphirhynchus is similar to Pseudoscaphirhynchus but differs in having a long, depressed, and completely armored caudal peduncle; it is endemic to the Mississippi basin of North America. The subfamily Husinae has two species: Huso dauricus (kaluga), which is native to the Amur River in Russia and China, and Huso huso (beluga), which occurs in the Black and Caspian seas and rarely in the Adriatic Sea. The beluga holds the distinction of being the largest freshwater fish in the world. Just how large is uncertain, as some records are in terms of length, others by weight. Hearsay has belugas reaching 6–8 m (20–26 ft) in length, but their weights are unknown. If such giants ever existed, it would have been well over 100 years ago. FishBase
Lake sturgeon (Acipenser fulvescens).
Acoela reports a 5-m (16.5-ft) 2072-kg (4558-lb) specimen estimated to be 118 years old. Sturgeons are sluggish, slow-growing, long-lived fishes found in Eurasia and North America. All spawn in freshwater, but several species spend most of their lives in the sea. Valued commercially for their flesh and as the source of caviar, sturgeons have long been seriously depleted in North America and most of Europe, but a substantial fishery is maintained in Western and Central Asia. See STURGEON. Family Polyodontidae (paddlefishes). Paddlefishes have a long, paddlelike snout. There are no large plates and scutes as in sturgeons, but patches of small ganoid scales; minute barbels on the underside of the snout; spiracles present; minute teeth; and greatly extended gill covers. There are several fossil genera known primarily from the Cretaceous Period in China and North America (Montana and Wyoming) and two living species, Psephurus gladius from China and Polyodon spathula from the United States. Polyodon differs from Psephurus in having long gill rakers that number in the hundreds and are used as a filter in feeding on plankton, in contrast to gill rakers that are shorter and fewer in number in Psephurus; and scales mostly on the caudal pecuncle and caudal fin in contrast to scales mostly on the trunk. Psephurus is the larger, reaching 3 m (10 ft) in total length. The largest Polyodon of recent record, weighing 65 kg (143 lb), was caught in the Missouri River in Montana in 1973. It is interesting to speculate on the utility of the paddle of the paddlefish. A common belief of nineteenth and early twentieth century fishers and professional naturalists was that the paddle is used to stir the bottom sediments to dislodge food, which then would be filtered from the water. Digging utility is not likely because the process would damage the numerous sensory and electroreceptors on the paddle. It is more likely that the paddle is used in detecting food. Another interesting hypothesis is that the paddle functions as a stabilizer to prevent dipping of the head that might occur as the fish cruises through water with its large mouth agape. Herbert Boschung; Reeve M. Bailey Bibliography. W. E. Bemis, E. K. Findeis, and L. Grande, An overview of Acipenseriformes, Environ. Biol. Fishes, 48:25–71, 1997; W. E. Bemis and B. Kynard, Sturgeon rivers: An introduction to Acipenseriformes biogeography and life history, Environ. Biol. Fishes, 48:167–183, 1997; V. J. Birstein and W. E. Bemis, How many species are there within the genus Acipenser?, Environ. Biol. Fishes, 48:157–163, 1997; V. J. Birstein, P. Doukakis, and R. DeSalle, Molecular phylogeny of Acipenseridae: Nonmonophyly of Scaphirhynchinae, Copeia, 2002(2):287–301, 2002; A. Choudhury and T. A. Dick, The historical biogeography of sturgeons (Osteichthyes: Acipenseridae): A synthesis of phylogenetics, paleontology and palaeogeography, J. Biogeog., 25:623–640, 1998; R. Froese and D. Pauly (eds.), FishBase, version 05/2005; L. Grande and
W. E. Bemis, Interrelationships of Acipenseriformes with comments on “Chondrostei,” pp. 85–115 in M. L. J. Stassny, L. R. Parenti, and G. D. Johnson (eds.), Interrelationships of Fishes, Academic Press, San Diego, 1996; L. Grande and W. E. Bemis, Osteology and phylogenetic relationships of fossil and recent paddlefishes (Polyodontidae) with comments on the interrelationships of Acipenseriformes, J. Vert. Paleontol., vol. 11, Mem. 1, 1991; J. S. Nelson, Fishes of the World, 4th ed., Wiley, New York, 2006.
Acoela An order of marine platyhelminthes of the class Turbellaria, generally regarded as primitive, and lacking protonephridia, oviducts, yolk glands, a permanent digestive cavity, and strictly delimited gonads. The nervous system is epidermal in the most primitive representatives and resembles that of cnidarians, but it is generally submuscular and consists of three to six pairs of longitudinal strands without anterior concentrations recognizable as a brain. The anterior or median ventral mouth opens into a simple pharynx which lacks elaborate musculature and leads directly to the syncytial vacuolated mass of endoderm forming the spongy intestinal tissue and lacking any true lumen. Eyes are often absent; when present they are of the usual paired turbellarian pigment-cup type, with one to a few pigment and retinal cells. A single (sometimes double) statocyst and sensory bristles (modified cilia) are generally present anteriorly, and clusters of frontal glands open by one or more pores at the anterior tip. The epidermis, typically, is uniformly ciliated, and locomotion is accomplished by means of rapid, smooth ciliary gliding or swimming. The acoels (see illus.) are mostly small (one to several millimeters in length), of simple contour, and plump and broad- to elongate-oval in form. Their
sensory bristles
frontal glands testes
statocyst
ovary
mouth egg sperm ducts seminal vesicle
Childia spinosa.
cuticular stylet of penis
73
74
Aconchulinida coloration is generally drab, or inconspicuous grayish-white, unless the body is colored brown or green by symbiotic algae. They are virtually worldwide in distribution, living beneath stones, among algae, on the bottom mud, or interstitially, from the littoral zone to deeper waters. A few are pelagic (in tropical and subtropical zones), and others have become symbiotic, living entocommensally in various invertebrates but mainly in echinoderms. Free-living species feed on diatoms, protozoans, and other microscopic organisms by partially extruding the digestive syncytium through the mouth and engulfing the food, ameboid pseudopodial fashion, in temporary vacuoles in the syncytium. Larger organisms (crustacean larvae, adult copepods, and the like) are seized by the anterior body (which may be expanded or curved dorsoventrally); the grip is facilitated by sticky mucus from the frontal glands, and the prey is thrust through the mouth into a larger temporary digestive cavity. Digestion is rapid and effected by acid and alkaline proteases, carbohydrases, and lipases acting in sequence and secreted from the syncytium. Entocommensal species feed and digest in the same manner, the presence of appropriate cosymbiotic organisms permitting retention of the basic nutritional physiology characteristic of their free-living relatives. The Acoela, as primitive Turbellaria, are of considerable interest phylogenetically since they show features reminiscent of the cnidarian grade of organization as well as others foreshadowing those of more elaborate turbellarians. In particular, their structure supports the hypothesis that ancient planula or planuloid larvae of hydrozoan cnidarians, by developing gonads and an oral aperture, gave rise to an acoeloid ancestral stock from which all modern flatworms, rhynchocoelans, and the coelomate Metazoa have subsequently evolved. See CNIDARIA; PLATYHELMINTHES; TURBELLARIA. J. B. Jennings
Aconchulinida An order of the subclass Filosia comprising a small group of naked amebas which form filamentous pseudopodia (filopodia). The genus Penardia, possi-
Aconchulinida. (a) Penardia cometa. (b) Penardia mutabilis.
bly the only valid one in the order, includes filopodiaforming amebas with variable shape, an outer zone of clear ectoplasm, and inner cytoplasm often containing many vacuoles and sometimes zoochlorellae (see illus.). Pseudopodia are sometimes extended primarily at the poles of the organism, but also may fuse laterally to form small sheets of clear cytoplasm. A single nucleus and one water-elimination vesicle are characteristic of these amebas. Size of described species ranges from less than 10 to about 400 micrometers. See PROTOZOA; RHIZOPODEA; SARCODINA; SARCOMASTIGOPHORA. Richard P. Hall
Acorales An order of monocotyledonous angiosperms composed of a single genus, Acorus (sweet flag or sweet calamus), with two species widespread in the Northern Hemisphere. Formerly, Acorus was included in the aroid family, Araceae, but several lines of evidence, including deoxyribonucleic acid (DNA) sequences, have firmly established it as the firstbranching lineage of the monocotyledons. These species are emergent aquatics with peculiar inflorescences, and they are unusual among the monocots in having characters otherwise confined to the magnoliid dicots, such as the anther formation and the presence of ethereal oils. These oils are the basis of the species’ frequent use as a medicine to treat toothache and dysentery, and as an ointment in religious ceremonies. See MAGNOLIOPHYTA. Mark W. Chase
Acoustic emission A method of nondestructive testing and materials characterization that uses mechanical waves moving through materials. It is similar to seismology, except in being concerned with the scale of engineering structures, such as aircraft, bridges, and chemical tanks. When a structure is subjected to external force (or stress), a defect (for example, a crack or welding flaw) on the structure is activated and enlarged dynamically, and thus generates waves, which spread through materials at a certain speed. Such waves, known as acoustic emission signals, are detected by sensors attached on the surfaces of the structure. Mechanical vibration due to acoustic emission signals is weak and requires high-sensitivity sensors and electronic amplification before it can be analyzed. See SEISMOLOGY. In nondestructive testing of structures, acoustic emission signals are typically evaluated in order to know if the failure of a structure is imminent; if cracks and other defects, presumed to be present in any structure, are active; the positions of such active defects; and whether a structure with such defects can be safely operated. In evaluating materials behavior and quality, acoustic emission is used to assess how a material (such as a steel or an aluminum alloy) responds to mechanical stress, that is, when and how it changes shape permanently and how it proceeds to
Acoustic emission 10 8
signal strength, V
6 4 2 0 −2 −4 −6 0
10
20 30 time, microseconds
40
Fig. 1. Waveform of rapidly varying acoustic emission signal from a brittle material.
ousness of the presumed defect is estimated, sometimes resulting in immediate shut-down of the entire plant. The sign of imminent failure is a rapid increase in the number of acoustic emission signals from a few concentrated sources. However, the value of the acoustic emission method lies in its pointing out potential trouble spots long before danger signs are generated. In this type of structural monitoring, the acoustic emission method can be applied without stopping normal operation (that is, without emptying the content of a tank or pipeline). Acoustic emission testing is applied to almost all tanks and vessels when new, and periodically during their service life. Other applications include detecting leakage of fluid and gas from pressurized tanks, pipes, and valves; monitoring key structural elements of bridges, dams, and tunnels; detecting damage in bearings and other rotating components, such as motors and compressors; detecting failure of cutting tools in automatic machining operation; and finding fissures and cracks in geothermal wells. In most 0.6 0.4 signal strength, V
eventual fracture; how an alloy withstands repeated application of stress (known as fatigue); the level of stress and corrosive environment that lead to failure of a material; and the types of microscopic failure processes that arise in a material under stress. See METAL, MECHANICAL PROPERTIES OF; PLASTIC DEFORMATION OF METAL; STRESS AND STRAIN. Methodology. Acoustic emission signals emanating from their sources contain information about the source, such as the direction and speed of crack opening. For example, the high-speed cracking of brittle materials (such as high-strength steels and ceramics) produces short, fast-varying acoustic emission signals (Fig. 1), which are typically plotted against time measured in microseconds. In contrast, slow-growing defects in plastics result in longer, slowly varying signals (Fig. 2), which are typically plotted on a time scale of milliseconds. Because of the distortion of waves during the propagation through a complex structure and detection by a sensor, however, much of the information is lost. Thus, the presence of detectable acoustic emission signals is the most important clue in assessing the integrity of the structure. By detecting one such signal at multiple sensor positions, the location of its source can be determined from the timing of signal arrivals. The basic principle of triangulation is the same as practiced in seismology, except that the differences in signal arrival times are of the order of microseconds to milliseconds. The speed of wave propagation is a material constant, determined by the stiffness and mass density of the propagating medium. Typically, the wave speed in solids is several kilometers per second. (By comparison, the sound speed in air is about 330 m or 1080 ft per second.) See SOUND. An example of a sensor arrangement is shown in Fig. 3 for a chemical tank. Here, the sensors are typically separated by several meters to 10 m (33 ft), and cover the entire surface of the tank with blocks of roughly triangular shape. The position of an acoustic emission source within each block is determined from the arrival times at the three surrounding sensors. A typical plot of acoustic emission source positions is shown in Fig. 4, indicating a concentration of sources. The detection of acoustic emission signals requires a sensor possessing high sensitivity. A typical sensor uses a piezoelectric ceramic element, which converts mechanical vibration into an electrical signal, which can be amplified 1000 to 10,000 times. Various electrical measurement techniques are used to characterize and analyze the signals received. It is common to obtain and record several features of acoustic emission signals. These form the basis of real-time analysis and decision-making. In laboratory studies, the entire waveforms are also recorded for detailed analysis after testing. See PIEZOELECTRICITY. Applications. By considering observed acoustic emission signal strength, correspondence to possible defective sites, and comparison to previous acoustic emission observations at the same position, the seri-
75
0.2 0 −0.2 −0.4 −0.6
0
5 10 time, milliseconds
15
Fig. 2. Waveform of slowly varying acoustic emission signal from a plastic.
50
Acoustic emission inside access
inside access (a) sensor
are serious defects and can be easily detected. One goal of materials engineering is to develop a strong and tough alloy with excellent fracture resistance, and acoustic emission provides valuable insights regarding the basic mechanical damage processes that are relevant to this endeavor. An example of damage accumulation is shown in Fig. 5. Here, stress is applied to a high-strength steel specimen. As the applied stress is increased beyond a limit, the number of acoustic emission signals starts to rise and continues to rise quickly. Each such signal is produced in this steel by the fracture of a very hard and brittle carbide particle embedded in the steel matrix. Near the final fracture point, some acoustic emission signals result from microscopic cracks. In this study, carbide particles are shown to be the main source of damage. Consequently, the strength of the steel is improved by 20% or more by reducing the size of carbide particles. See STEEL. In addition to the observation of deformation and fracture, acoustic emission is utilized in various areas of materials studies. These include abrasion, casting,
10 sensor
inside access
sensor 12
(b) Fig. 3. Sensor arrangement for a chemical tank. (a) Side view. (b) Top view.
sensor
9
Fig. 4. Typical plot of acoustic emission source positions, indicating a concentration of sources.
fracture
2 cumulative number of acoustic emission signals
8000 6000
1
applied stress
4000 2000
0
cumulative number of acoustic emission signals
instances, the phenomena of acoustic emission generation are complex, and the evaluation of the results is largely learned from experience. Acoustic emission methods are also valuable in the laboratory. Acoustic emission signals are emitted when materials are permanently deformed (stretched, pressed, bent, or twisted). Materials can exhibit several characteristic acoustic emission behaviors. One type of material emits many acoustic emission signals when the deformation begins, followed by a decreasing trend as more deformation occurs. Another type is characterized by the near absence of acoustic emission. This behavior is found when the material has been previously deformed heavily. These changes reflect the variation in the internal composition of materials, and acoustic emission provides key information in understanding deformation characteristics. Acoustic emission signals are also produced when materials are stressed to a level nearing fracture. Material fracture results from the accumulation of microscopic internal damage. Such damage develops either gradually or rapidly. While acoustic emission cannot detect slow damage formation (such as forming voids), rapid damage processes that form cracks
applied stress, gigapascals
76
0 0
1 2 deformation, mm
Fig. 5. Behavior of a high-strength tool steel specimen as stress is applied.
Acoustic impedance film coating and plating, corrosion, curing of plastics, friction, magnetization, oxidation, and hardening of steels. See NONDESTRUCTIVE EVALUATION. Kanji Ono Bibliography. Acoustic emission, in Metals Handbook, 9th ed., vol. 17: Nondestructive Testing and Quality Control, pp. 278–294, ASM International, Materials Park, OH, 1989; R. K. Miller and P. McIntire (eds.), Acoustic emission testing, Nondestructive Testing Handbook, 2d ed., vol. 5, American Society for Nondestructive Testing, Columbus, OH, 1987; K. Ono (ed.), J. Acous. Emission, 1982– ; K. Ono (comp.), Progress in Acoustic Emission IX, 1998 [Proceedings of the 14th International Acoustic Emission Symposium: Transitions in AE for the 21st century]; I. G. Scott, Basic Acoustic Emission, Gordon and Breach, New York, 1991.
Acoustic impedance The ratio of the complex amplitude of spatially averaged acoustic pressure over a surface to the complex amplitude of volume velocity on that surface. The unit is the newton-second/meter5, or the mks acoustic ohm. In the cgs system the unit is the dynesecond/centimeter5. The volume velocity is the integral over the surface of the particle velocity normal to the surface. Acoustic impedance measures the degree to which an acoustic medium impedes the fluid motion created by an imposed pressure. It was originally defined by analogy to electrical systems. See ELECTRICAL IMPEDANCE; SOUND PRESSURE; UNITS OF MEASUREMENT. In cases where the acoustic pressure p and volume velocity Q vary sinusoidally with time t, it is convenient to express them (as with quantities in alternating-current circuits) as the real parts of complex numbers that are rotating in the complex plane. ˜ are complex ˜ and Q This is done in Eqs. (1), where p ˜ exp(iωt)} p(t) = Re { p ˜ exp(iωt)} Q(t) = Re {Q
(1)
√ numbers, Re means “real part of,” i = −1, and ω is the angular frequency. The acoustic impedance is given by Eq. (2). It is generally complex-valued beZ=
˜ p ˜ Q
(2)
˜ are generally complex˜ and Q cause the amplitudes p valued. The phase of Z represents the phase difference between pressure and volume velocity. For a fixed pressure, a large volume velocity is associated with a small magnitude of acoustic impedance. See ALTERNATING-CURRENT CIRCUIT THEORY; COMPLEX NUMBERS AND COMPLEX VARIABLES. Specific acoustic impedance z is defined as the acoustic impedance multiplied by the area, so given by Eq. (3), where V˜ is the complex amplitude of the z = AZ =
˜ p ˜ V
(3)
K M
f (t)
C V (t)
Mechanical system used to interpret acoustic impedance.
average velocity normal to the surface. The unit is the newton-second/meter3, or the mks rayl. In the cgs system the unit is the dyne-second/centimeter3, or the rayl. The characteristic acoustic impedance is the specific acoustic impedance for a plane wave propagating through an acoustic medium and is equal to ρc, where ρ is the mass density and c is the sound speed. The characteristic impedance of water is approximately 1.5 × 106 mks rayls, and the characteristic impedance of air at 0◦C (32◦F) is approximately 400 mks rayls. Interpretation of real and imaginary parts. Acoustic impedance is often expressed as the sum of real and imaginary parts, as in Eq. (4), where R is the acousZ = R + iX
(4)
tic resistance and X is the acoustic reactance. These quantities may be interpreted by analogy to the mechanical system shown in the illustration, where the mechanical impedance is defined by Eq. (5). A Zmech =
f˜ K = C + i ωM − ω V˜
(5)
comparison of Eqs. (4) and (5) indicates that the resistance is analogous to a mechanical dashpot (a damping element in which the applied force is proportional to velocity) and the reactance is analogous to either a mass or a spring, depending on its sign. A more general interpretation of the resistance comes from Eq. (6) for the time-averaged power flow˜ 2 Pave = 1/2 R|Q|
(6)
ing through the surface, which indicates that the power flow is proportional to resistance for a fixed volume velocity. The reactance is associated with an energy storage. Circular aperture in planar baffle. One widely used example of acoustic impedance arises when the surface is chosen as a circular aperture in a planar rigid baffle of infinite extent with an acoustic fluid on one side. The velocity normal to the aperture is assumed constant over the aperture. The acoustic fluid has mass density ρ and sound speed c. When the radius of the circle is much less than an acoustic wavelength,
77
78
Acoustic interferometer the resistance and reactance are given by Eqs. (7), Rp = ρc
2
1 (ka) 2 πa2
8 ka Xp = ρc 3π πa2
(ka)4 y*) at temperature T * (T * > T0). The column (section 5) behind the rear zones is equilibrated with the feed gas at feed conditions. (2) In type II behavior, a pure heat-transfer zone is formed, followed by a pair of mass- and heattransfer zones as shown in Fig. 3. The adsorbate is absent in sections 1 to 3 in this case. The column
Adsorption operations
(type I)
(8a)
n0 cs > y0 cg
(type II)
(8b)
4
3
2
1
y* T*
flow direction 0
temperature
mole fraction
5
T0
distance in column
Fig. 2. Composition and temperature profiles during adsorption of type I systems in packed columns. Arrows at the curves indicate the relevant axes.
4
3
2
1
y0 T*
temperature
mole fraction
5
y/yo
component 1 flow direction 0 distance in column
Fig. 4. Roll-over effect showing displacement of more weakly adsorbed component 2 by component 1 during isothermal adsorption from an inert carrier gas.
y* TH y0 flow direction 0
conditions and cs [ J/(g − K)] and cg [ J/(mol − K)] are, respectively, the heat capacities of the adsorbent and the gas mixture. Type I behavior is common for adsorption of bulk adsorbate systems, while type II is favored when the adsorbate is dilute (y0 1) and strongly adsorbed. Multiple transfer zones and equilibrium sections are formed in multicomponent adsorption, which is also characterized by roll-over effects, where the more strongly adsorbed species displaces the weaker
y0
component 2
flow direction
T0
0 distance in column
Fig. 3. Composition and temperature profiles during adsorption of type II systems in packed columns. Arrows at the curves indicate the relevant axes.
temperature
n0 cs < y0 cg
1.0
mole fraction
behind the rear zones remains equilibrated with the feed gas at feed conditions. In both cases, the zones propagate through the column after their formation as more feed gas is introduced. The front zones move faster than the rear zones. The shapes and sizes of the zones can be measured by monitoring the y(t) and T(t) profiles in the effluent gas from the column. The effluent profiles are known as breakthrough curves, which are mirror images of the profiles within the column. The zones may expand (proportionate pattern) or retain their shape (constant pattern) as they move through the column. Constant-pattern front zones for type I systems and constant-pattern rear zones for type II systems are usually formed in a long column when the isotherm of the adsorbate is sufficiently concave toward the pressure axis. Approximate criteria for the formation of type I and type II systems are given in Eqs. (8), where n0 is the adsorbate capacity at feed
distance in column
T0
Fig. 5. Composition and temperature profiles during thermal desorption in packed columns. Arrows at the curves indicate the relevant axes.
ones as the zones propagate. Figure 4 shows, for example, y profiles for an isothermal column adsorbing two substances from an inert gas; component 1 is more strongly adsorbed. The general behavior of a column can be estimated by the simultaneous solution of the partial differential equations describing the mass and heat balances for each adsorbate species present in the column, in conjunction with their kinetic and equilibrium adsorption properties. Numerical solutions requiring considerable computation time are generally needed. Other factors influencing the dynamics are gas maldistribution, channeling, and nonadiabatic columns. Desorption. Desorption of adsorbates from a column is usually accomplished by heating the column with a hot, weakly adsorbed gas; lowering the column pressure; purging the column with a weakly adsorbed gas; or combinations of these methods. Like adsorption, well-defined transfer zones and equilibrium sections can be formed during desorption by the methods of heating and purging as shown in Figs. 5 and 6, respectively, for the case of a single adsorbate. Two pairs of mass- and heat-transfer zones are formed in both cases. The sections ahead of the front zones remain equilibrated at the initial column conditions. The middle equilibrium sections, however, attain temperatures (T *) and compositions (y*) greater or less than T0 and y0 depending on whether the desorption is by the heating or the purging method, respectively. The section behind the rear transfer zones is equilibrated at the conditions of the purge gas.
167
yo flow direction
To T*
y*
temperature
Aerial photography
mole fraction
168
0 distance in column Fig. 6. Composition and temperature profiles during purge desorption in packed columns. Arrows at the curves indicate the relevant axes.
Adsorbents. Microporous adsorbents like zeolites, activated carbons, silica gels, and aluminas are commonly employed in industrial gas separations. These solids exhibit a wide spectrum of pore structures, surface polarity, and chemistry which makes them specifically selective for separation of many different gas mixtures. Separation is normally based on the equilibrium selectivity. However, zeolites and carbon molecular sieves can also separate gases based on molecular shape and size factors which influence the rate of adsorption. See ACTIVATED CARBON; MOLECULAR SIEVE; ZEOLITE. Industrial applications. By far the most frequent industrial applications of gas adsorption have been the drying of gases, solvent vapor recovery, and removal of impurities or pollutants. The adsorbates in these cases are present in dilute quantities. These separations use a thermal-swing adsorption process whereby the adsorption is carried out at a nearambient temperature followed by thermal regeneration using a portion of the cleaned gas or steam. The adsorbent is then cooled and reused. Many variations of this theme are practiced to reduce the energy consumption and the adsorbent inventory. Pressure-swing adsorption processes are employed for separation of bulk gas mixtures. The adsorption is carried out at an elevated pressure level to give a product stream enriched in the more weakly adsorbed component. After the column is saturated with the strongly adsorbed component, it is regenerated by depressurization and purging with a portion of the product gas. The cycle is repeated after raising the pressure to the adsorption level. Numerous combinations of these cyclic steps in conjunction with internal recycling of gases have been patented. Key examples of pressure-swing adsorption processes are production of enriched oxygen and nitrogen from air; production of ultrapure hydrogen from various hydrogen-containing streams such as steam-methane reformer off-gas; and separation of normal from branched-chain paraffins. Both thermalswing and pressure-swing adsorption processes use multiple adsorbent columns to maintain continuity, so that when one column is undergoing adsorption the others are in various stages of regeneration modes. The thermal-swing adsorption processes typically use long cycle times (hours) in contrast to the rapid cycle times (minutes) for pressure-swing adsorption processes. A. L. Myers; S. Sircar
Bibliography. S. J. Gregg and K. S. W. Sing, Adsorption, Surface Area and Porosity, 1982; A. L. Myers and G. Belfort (eds.), Fundamentals of Adsorption, 1984; R. H. Perry and D. Green (eds.), Perry’s Chemical Engineers’ Handbook, 7th ed., 1997; D. M. Ruthven, Principles of Adsorption and Adsorption Processes, 1984; D. M. Young and A. D. Crowell, Physical Adsorption of Gases, 1962.
Aerial photography A photograph of a portion of the Earth’s surface taken from an aircraft or from a satellite. Most often, these photographs are taken sequentially and overlap each other, to give complete coverage of the area of interest. Thus they may be viewed stereoscopically to give a three-dimensional view of the Earth’s surface (Fig. 1). Although the camera in the aircraft may be pointed obliquely to the side of the line of flight, by far the most common type of aerial photograph is taken with the camera pointed vertically downward beneath the plane. See STEREOSCOPY. The entire United States has been photographed, generally more than once, and most of the rest of the world has also been covered. In the United States, the federal government is the major holder and seller of aerial photographs, although private companies also make and sell them. Federal agencies such as the Geological Survey, the Department of Agriculture, NASA, and others routinely use aerial photographs. Aerial photographs may be made with a variety of films and filters. By far the most common is blackand-white panchromatic film used with a yellow filter. Other films are black-and-white infrared, color, and color infrared. The color and color infrared films are usually flown by companies on contract with the user and are rarely available from the government. The scales of aerial photographs are generally between 1 to 2500 and 1 to 80,000, and the majority have been taken at a scale of about 1 to 20,000. On a 1-to-20,000 photograph, 1 mi on the ground is about 3 in. on the photograph (1 km = 5 cm). As the most frequently available picture size is 9 in. by 9 in. (23 cm by 23 cm), most available federal aerial photographs cover about 9 mi2 (23 km2) each. Photogrammetry. Aerial photographs have two main uses, the making of planimetric and topographic maps, and interpretation or data-gathering in a variety of specialized fields. All modern topographic (contour) maps are made from stereoscopic aerial photographs, usually black-and-white panchromatic vertical photographs. The science of making accurate measurements and maps from aerial photographs is called photogrammetry. Because vertical aerial photographs are central perspective views of the Earth (seen through a single small camera lens), and because the camera may not be precisely vertical when the picture is taken, the individual aerial photograph is not an exact map or geometric representation of the Earth’s surface. The geometric discrepancy between the distribution of features seen on the photograph and their actual distribution on the
Aerial photography
Fig. 1. Stereogram of sedimentary rocks outcropping around eroded structural dome in Wyoming. Such photographs supply much of the information needed in preparing geologic maps. Structures of this type are commonly important in localizing petroleum accumulations. To obtain a three-dimensional image, a standard lens-type stereoscope may be used. (USGS)
Earth’s surface depends on several variables, the most important of which is the relief of the surface. The error becomes greater as relief increases. Photogrammetry deals with these errors and allows highly accurate and complete map production. Topographic maps are made from one or more stereoscopic pairs of aerial photographs by means of optical-mechanical plotting machines or mathematical methods (analytical photogrammetry). In addition to topography, many other features such as vegetation and structures made by people can be transferred from the photographs to the maps. There is not now available any method of map production which can compete in versatility, speed, and completeness with aerial photographs. See CARTOGRAPHY. Photo interpretation. The other major use of aerial photographs is photo interpretation, utilizing all types of film. Photo interpretation is the attempt to extract information about the Earth’s surface (and sometimes, subsurface) from aerial photographs. It is usually carried out by specialists in a subject area. Some users of photo interpretation methods are geologists, foresters, hydrologists, agricultural specialists, soil scientists, land-use planners, environmentalists, military intelligence personnel, engineers, and archeologists. The systematic study of aerial photographs, particularly when they are viewed stereoscopically, gives these specialists the ability to gather rapidly, record, map, and interpret a great deal of
information. In addition, an aerial photograph is a record at a moment in time; photographs taken at intervals give an excellent record of the change with time in surface features of the Earth. In some specialties this is of vital importance. The use of different types of film, such as black-and-white infrared or color infrared, allows the interpreter to see the Earth’s surface by reflected light energy beyond the range of the human eye. Thus features normally difficult to see may be made apparent. Applicability. Some examples of the use of aerial photographs may make clear their value. Geologists use the photographs primarily for geologic mapping. It is often possible to see different types of rocks and trace their distribution. This, combined with fieldwork, enormously increases the efficiency and completeness of geologic mapping. But, in addition, the synoptic view of the Earth’s surface given by the photograph enables the geologist to map fractures in the Earth’s surface, subtle tonal variation in the soil, and other features difficult or impossible to see on the ground. The land-use planner obtains a synoptic view of the surface and of human activities. This, coupled with aerial photographs taken at various times, gives a spatial and temporal database vital to planning. The forester has learned to recognize various tree types and calculate tree sizes, thus enabling rapid and accurate forest inventories to be done. Importantly, various types of film allow rapid recognition of diseased trees and mapping of the
169
170
Aerodynamic force siderable importance in the interpretation of satellite imagery, such as that transmitted from Landsat. In this case the original output is not a photograph but is digital and may be treated mathematically before being made into a photographic image. The field of remote sensing is much concerned with these techniques. Aerial photographs, as well as satellite imagery and airborne radar imagery, are remotesensed images which, taken as a group, provide a view of the Earth’s surface that was never before available (Fig. 2). See REMOTE SENSING. L. H. Lattman Bibliography. J. Ciciarelli, A Practical Guide to Aerial Photography, 1991; R. Graham, Digital Aerial Survey: Theory and Practice, 2002; H. Lloyd, Aerial Photography: Professional Techniques and Applications, 1990; E. M. Mikhail et al., Introduction to Modern Photogrammetry, 2001; D. P. Paine and J. D. Kiser, Aerial Photography and Image, 2003.
Aerodynamic force
Fig. 2. Aerial view over Sinai Peninsula from Gemini II, looking northeast. Eastern desert of Egypt is in foreground, with Israel and Jordan in background. Note detail in which topographic and drainage features are shown. (NASA)
extent of the disease. Agricultural specialists can recognize and map various crops, thus making crop inventories in much less time and with much less expense than could be done purely by ground methods. Disease, drought, and other crop stresses can be recognized and mapped. Archeologists can see ancient occupation sites, travel routes, and other features on aerial photographs which are very difficult to see on the ground. Environmental studies are greatly aided by the enlargement of the aerial photograph, and environmental monitoring gains much from the study of changes seen on aerial photographs taken at different times. The use of aerial photographs for military intelligence gathering, bomb damage assessment, and other uses is well known. This short, incomplete list of the wide applicability of photo interpretation to many fields of study indicates the great value of aerial photographs. See AGRICULTURE; LANDUSE PLANNING; TOPOGRAPHIC SURVEYING AND MAPPING; VEGETATION AND ECOSYSTEM MAPPING. Remote sensing. Photo interpretation depends on two things: the appearance and distribution (geometry) of familiar objects, and their spectral reflectance characteristics. The latter category is of increasing importance and is based on the fact that objects on the Earth’s surface reflect different amounts of energy (light) in different parts of the electromagnetic spectrum. Thus by careful selection of films and filters with different light sensitivity in different wavelengths of the spectrum, individual “target” objects may be emphasized. This technique is of con-
The force exerted on a body whenever there is a relative velocity between the body and the air. There are only two basic sources of aerodynamic force: the pressure distribution and the frictional shear stress distribution exerted by the airflow on the body surface. The pressure exerted by the air at a point on the surface acts perpendicular to the surface at that point; and the shear stress, which is due to the frictional action of the air rubbing against the surface, acts tangentially to the surface at that point. The pressure and shear stress act at each point on the body surface, and hence the distribution of pressure and shear stress represent a distributed load over the surface. The net aerodynamic force on the body is due to the net imbalance between these distributed loads as they are summed (integrated) over the entire surface. See BOUNDARY LAYER FLOW; FLUID FLOW; WIND STRESS. Lift and drag definitions. For purposes of discussion, it is convenient to consider the aerodynamic force on an airfoil (Fig. 1). The net resultant aerodynamic force R acting through the center of pressure R
L
center of pressure
velocity, V ∞ density, ρ∞
D
α
c
Fig. 1. Resultant aerodynamic force (R), and its resolution into lift (L) and drag (D) components.
Aerodynamic force
2 SCL L = 1/2 ρ∞ V∞
(1)
2 /2 ρ∞ V∞ SCD
(2)
D=
1
the lift coefficient and CD as the drag coefficient. These coefficients vary with angle of attack. At higher speeds, especially near and above the speed of sound, CL and CD also become functions of the freestream Mach number M∞, defined as M∞ = V∞/a∞, where a∞ is the speed of sound in the free stream. Also, in flows where friction is an important aspect of the aerodynamic force, CL and CD become functions of the Reynolds number, defined as Re = ρ ∞V∞c/µ∞, where c is some characteristic length of the body (such as the chord length c; Fig. 1), and µ∞ is the coefficient of viscosity. Therefore, in general, the functional dependence of the lift and drag coefficients is given by Eqs. (3) and (4). CL = f1 (α, M∞ , Re)
(3)
CD = f2 (α, M∞ , Re)
(4)
See MACH NUMBER; REYNOLDS NUMBER. Separation and stall. The lift varies linearly with angle of attack, up to a point where the curve that graphs the lift as a function of angle of attack bends over (Fig. 2). Beyond this point the lift decreases with increasing α. In this region, the wing is said to be stalled. On the linear portion of the lift curve, the flow over the airfoil is smooth and attached (Fig. 2). In contrast, in the stall region the flow has separated from the top surface of the wing, creating a type of
separated flow
CL
CL, CD, L/D
on the airfoil represents mechanically the same effect as that due to the actual pressure and shear stress loads distributed over the body surface. The velocity of the airflow V∞ is called the free-stream velocity or the free-stream relative wind. By definition, the component of R perpendicular to the relative wind is the lift, L, and the component of R parallel to the relative wind is the drag D. See AIRFOIL. Force coefficients. The orientation of the body with respect to the direction of the free stream is given by the angle of attack, α. The magnitude of the aerodynamic force R is governed by the density ρ ∞ and velocity of the free stream, the size of the body, and the angle of attack. An index of the size of the body is given by some characteristic area S of the body. If the body is a winglike shape, a suitable choice for S is the planform area (the area seen when the wing is viewed directly from above). If the body is shaped like a projectile, a suitable choice for S is the maximum cross-sectional area of the body. Since S is used simply as an index of the size of the body, the area chosen is somewhat arbitrary. The aerodynamic force on a given body shape moving through the air at speeds considerably less than the speed of sound is directly proportional to the density, the area, and the square of the velocity; hence, the same proportionality holds for lift and drag. The proportionality constants are denoted by CL/2 and CD/2, respectively, so that the lift and drag are given by Eqs. (1) and (2). Here, CL is defined as
L /D
attached flow
CD α Fig. 2. Typical variation of lift coefficient (CL), drag coefficient (CD), and lift-to-drag ratio (L/D) for an airfoil as functions of angle of attack (α). Insets show physical nature of the flow field over the wing at the indicated points on the lift curve.
Maximum lift-to-drag ratios of representative aircraft Aircraft
(L/D)max
Wright flyer Sopwith Camel Ryan NYP (Spirit of St. Louis) Douglas DC-3 Boeing B-52 F-4 fighter
5.7 7.7 10.1 14.7 21.5 8.6
slowly recirculating dead-air region, which decreases the lift and substantially increases the drag. An important measure of aerodynamic efficiency is the ratio of lift to drag, L/D. The higher the value of L/D, the more efficient is the lifting action of the body. The value of L/D reaches a maximum, denoted by (L/D)max (see table), at a relatively low angle of attack (Fig. 2). See FLIGHT CHARACTERISTICS. Drag components. The total aerodynamic drag on a body is often dissected into contributing parts. Skin friction drag. This is due to the net integrated effect of the surface shear stress exerted over the body. Pressure drag due to flow separation. When the flow separates from the body, the change in the surface pressure distribution in the dead-air separated region causes an increase in drag. Although the aerodynamic mechanism that causes the flow to separate from the surface involves the effect of friction in the flow, this drag is purely a pressure effect. This drag is sometimes called form drag. Parasite and profile drag. The sum of skin-friction drag and pressure drag due to flow separation is denoted for a wing as profile drag, and for a complete airplane as parasite drag. Induced drag. This is associated with the vortices that are produced at the tips of wings and trail downstream from the wings (Fig. 3). A wing produces lift because the pressure on its bottom surface is
171
172
Aerodynamic sound vortex low pressure high pressure (a) freestream velocity
having to do with the shape of the planform of the wing; values of e typically range from 0.85 to 1.0. As seen in Eq. (6), induced drag increases as the square of the lift. See ASPECT RATIO. For a complete airplane at subsonic speed, the total drag coefficient is the sum of the parasite drag coefficient CDP and induced drag coefficient, as in Eq. (7). Since the pressure drag on an airplane deCD = CDP +
(b)
CD = CD,0 +
Fig. 3. Origin of wing-tip vortices on a finite wing as shown in (a) front view and (b) oblique view of wing.
higher than that on the top. At a wing tip, the higherpressure air curls around the tip from bottom to top, setting up a circulatory motion which is superimposed on the main flow velocity and creates vortices which trail downstream from each tip. In turn, the presence of these vortices changes the pressure distribution over the wing surface in comparison to that which would exist if the wing had no tips (that is, if the wing were infinite in span). This change in the pressure distribution is in such a direction that it results in an increase in drag. This increment in drag is called induced drag; it is purely a pressure effect. Since the wing-tip vortices, which cause induced drag, in turn result from the pressure differential on the wing surfaces, and this pressure differential is also responsible for the creation of lift, induced drag is frequently called drag due to lift. See VORTEX. Induced drag is the penalty paid for the production of lift. For an aircraft in steady, level flight, the lift produced must equal the weight of the aircraft. The induced drag caused by this lift must be overcome by increased thrust from the engines. Wave drag. This is associated with the presence of shock waves on the aerodynamic vehicle. It is a pressure drag that is caused by severe increases in pressure that occur across shock waves. See AERODYNAMIC WAVE DRAG. Drag coefficient expressions. In essence, the total drag on a body can be visualized as the sum of parasite drag, induced drag, and wave drag. For a wing at subsonic speeds, the total drag coefficient CD is expressed as the sum of the profile drag coefficient cd, p and induced drag coefficient CDi, as in Eq. (5), where the induced drag coefficient is given by Eq. (6). Here,
CDi =
CL2 πeAR
(7)
pends on angle of attack, the total drag for an airplane is frequently expressed as a sum of the zero-lift drag and the drag due to lift. In coefficient form, the total airplane drag coefficient is given by Eq. (8). Here,
wing-tip vortices
CD = cd,p + CDi
CL2 πeAR
(5) (6)
AR is the wing aspect ratio, defined as b2/S, where b is the span of the wing (distance from one wing tip to the other), and e is a span effectiveness factor
CL2 πe1 AR
(8)
CD,0 is the parasite drag coefficient at zero lift, and C2L /πe1AR is the drag coefficient due to lift, which is physically the sum of the induced drag coefficient plus the extra increment in parasite drag coefficient when the airplane is pitched to an angle of attack to produce lift. This extra parasite drag is reflected in the factor e1, called Oswald’s efficiency factor, whose value is typically of the order of 0.55–0.7. Equation (8) is called the drag polar for the airplane, and contains all the practical aerodynamic information needed to analyze the airplane’s aerodynamic performance. John D. Anderson, Jr. Bibliography. J. D. Anderson, Jr., Fundamentals of Aerodynamics, 4th ed., 2005; J. D. Anderson, Jr., Introduction to Flight, 5th ed., 2005; J. J. Bertin, Aerodynamics for Engineers, 4th ed., 2002; S. F. Hoerner, Fluid-Dynamic Drag, 1965; S. F. Hoerner and H. V. Bortst, Fluid-Dynamic Lift, 2d ed., 1985.
Aerodynamic sound Sound that is generated by the unsteady motion of a gas and its interaction with surrounding surfaces. Aerodynamic sound or noise may be pleasant, such as the sound generated by a flute, or unpleasant, such as the noise of an aircraft on landing or takeoff, or the impulsive noise of a helicopter in descending or forward flight. The typical human ear can detect a great range of sounds. A logarithmic scale, called the decibel (dB) scale, is used to describe the sound pressure level (SPL). The sound pressure level is given by the equation below, where prms = root-mean-square prms dB SPL = 20 log10 pref pressure level of the sound, and reference pressure pref = 2 × 10−5 N/m2. Zero decibels corresponds to the quietest whisper, whereas the threshold of pain occurs at 130–140 dB. See LOUDNESS; SOUND PRESSURE. Basic principles. Sources of aerodynamic sound may be classified according to their multipole order. Sources associated with unsteady mass addition to
Aerodynamic sound the gas are called monopoles. These could be caused by the unsteady mass flux in a jet exhaust or the pulsation of a body. The sound power radiated by monopoles scales with the fourth power of a characteristic source velocity. Sources related to unsteady forces acting on the gas are called dipoles. The singing in telephone wires is related to the nearly periodic lift variations caused by vortex shedding from the wires. Such sources, called dipoles, generate a sound power that scales with the sixth power of the characteristic source velocity. See KA´ RMA´ N VORTEX STREET. A turbulent fluid undergoes local compression and extension as well as shearing. These events are nearly random on the smallest scales of turbulent motion but may have more organization at the largest scales. The earliest theories of aerodynamic noise, called acoustic analogies, related these unsteady stresses in the fluid to the noise that they would generate in a uniform ambient medium with a constant speed of sound. Such sources are called quadrupoles. The sound power that they radiate scales with the eighth power of the characteristic source velocity. See TURBULENT FLOW. Subsequent extensions of these theories allowed for the motion of these sources relative to a listener. This may result in a Doppler shift in frequency and a convective amplification of the sound if the sources are moving toward the listener. In addition, if the sources are embedded in a sheared flow, such as the exhaust plume of a jet engine, the sound is refracted away from the jet downstream axis. As sound propagates away from the source region, it experiences attenuation due to spherical spreading and real-gas and relaxational effects. The latter effects are usually important only for high-amplitude sound or sound propagation over large distances. See ATMOSPHERIC ACOUSTICS; DOPPLER EFFECT; SOUND ABSORPTION. Aircraft noise. The main sources of aircraft noise are the jet engine and the airframe. Propellerdriven aircraft and helicopters have additional noise sources. Engine noise. A modern commercial jet engine has a high bypass ratio, that is, the ratio of the mass flow rates through the fan and core (Fig. 1). An alternative design mixes the fan and core streams internally and has a single exhaust nozzle. The sound power generated by the jet exhaust plume is proportional to the eighth power of the exit velocity and the square
fan
exit guide vanes
jet exhaust jet nozzle nacelle Fig. 1. Separate-flow high-bypass-ratio jet engine.
of the exit area. The high bypass ratio reduces the amount of thrust generated by the high-speed flow and replaces it with thrust delivered by the relatively low-speed fan flow. This reduces the overall radiated noise power. For jets with a subsonic exit velocity, turbulent mixing generates stresses in the exhaust that in turn generate the radiated noise. The directivity of the jet noise peaks at a relatively small angle to the jet downstream direction. This is a result of the convective amplification and refraction effects. At high jet exit velocities, the turbulence may convect supersonically with respect to the ambient speed of sound. This results in highly directional and efficient noise radiation called Mach wave radiation. These are all examples of jet mixing noise. See AIRCRAFT ENGINE PERFORMANCE. If the jet is operating off-design, that is, the jet exit static pressure is different from the ambient pressure, a shock cell structure forms in the jet exhaust. Shock cells are diamond-shaped regions of alternating high and low pressure in the jet exhaust through which the jet exit pressure adjusts to the ambient pressure. The interaction between the turbulence in the jet and the shock cell structure generates broadband shock-associated noise and screech. The former noise radiates primarily in the upstream direction, is broadband, and peaks at frequencies higher than the jet mixing noise. Screech is a feedback phenomenon in which sound generated by the turbulence interacting with the shock cells propagates upstream and triggers turbulence at the jet exit. This convects downstream, interacting again with the shock cells and completing the cycle. The feedback cycle generates an intense tone known as jet screech. Noise is generated by the fan in the high-bypassratio engine due to interactions between the wakes shed by the fan blades and the exit guide vanes. This interaction generates spinning acoustic modes that may propagate in the fan duct, either upstream toward the engine inlet or downstream through the fan exhaust duct. (Spinning acoustic modes are the natural acoustic waveforms in a circular duct. They consist of pressure disturbances with wave fronts that follow helical paths.) The noise consists of a tonal component proportional to the fan’s blade passage frequency and a broadband component associated with the turbulent wake interactions as well as turbulence ingested into the fan duct. The turbulent interactions cause unsteady loading on the fan blades and the exit or outlet guide vanes. The reaction forces acting on the fluid generate noise. Acoustic liners, usually made of perforated plates with a honeycomb backing, can attenuate the duct modes. In addition, the choice of the relative number of fan blades to exit guide vanes, as well as the tip speed of the fan, affects the noise generation efficiency. See AIRCRAFT ENGINE PERFORMANCE; JET PROPULSION; TURBOFAN; WAKE FLOW. Airframe noise. Airframe noise is generated by the nonpropulsive components of an aircraft. With continued reduction in engine noise, airframe noise can be a significant component of the overall aircraft noise, especially during landing. The major sources
173
174
Aerodynamic wave drag
slats
nose gear
flaps main landing gear
Fig. 2. Sources of airframe noise.
of airframe noise are associated with the high-lift devices, namely the trailing-edge flaps and leading-edge slats, and the landing gear—both the nose and main gear (Fig. 2). See FLIGHT CONTROLS; LANDING GEAR. Flap systems may consist of multiple elements. Also, they may not extend over the entire wingspan. Gaps or gates may be used to limit engine exhaust impingement on the flaps. Vortices are shed from the flap side edges. The strength of the vortices depends on the loading distribution on the flap. The vortex moves over the flap side edge, inducing unsteady surface pressures on the flap, resulting in dipole noise radiation. In addition, the turbulence in the vortex generates noise that can scatter from the side edge corners or flap trailing edge. Since the loading noise scales with the sixth power of the local Mach number, and corners and edges scatter noise that scales with the fifth power, the latter mechanism can dominate at low Mach number. See MACH NUMBER; VORTEX. Noise is generated at the trailing edge of the leading-edge slat. This can be tonal and is believed to be associated with vortex shedding from the slat trailing edge. This vortex shedding may be reinforced by a resonant feedback in which the radiated noise reflects from the main element of the wing. The vortex shedding depends on the slat trailing-edge thickness as well as the boundary layer thicknesses on the upper and lower surfaces of the slat. The landing gear acts as a bluff body and generates noise associated with unsteady loading in a similar manner to the vortex-shedding noise of a circular cylinder. If any sharp edges are in the vicinity of the landing gear, such as a door that opens with the gear, acoustic scattering may occur, which increases the efficiency of the noise radiation. In addition, smallscale fittings, such as hydraulic lines attached to the gear, generate noise at a higher frequency, due to their smaller scale, than that generated by the wheels and the main gear components. Propeller and rotorcraft noise. The noise generated by propellers contains a broadband and a harmonic component. The harmonic component occurs at multiples of the blade passage frequency. It consists of thickness noise and loading noise. Thickness noise is caused by the displacement of the air by the rotating propeller blade. It is a monopolelike source acting at the blade surface. Loading noise is related to the lift and drag forces acting on the blade. The source is of dipole type. For a propeller not in for-
ward motion, the noise behind the plane of the propeller blades is greater than that ahead of the propeller. Forward flight causes convective amplification so that noise levels are greater when the aircraft is approaching a stationary observer. As is the case for fan noise, broadband propeller noise is associated with turbulence ingestion and blade self-noise. The latter sources include the surface pressure fluctuations caused by the turbulent boundary layer on the blade and vortex shedding from the blade trailing edge. If the propeller tip speed is supersonic, sound generated at several points on the blade surface can reach a listener at the same time. The reason is that the source region is moving faster relative to the listener than the local speed of sound. The listener hears an intense, crackling sound during a short fraction of the propeller’s rotation period. See PROPELLER (AIRCRAFT). The same harmonic and broadband sources arise in helicopter noise generation. The noise levels for a helicopter in hover are greater below the plane of the rotor. In addition, in descending flight the rotor blades may encounter tip vortices shed from the passage of previous blades. This gives a transient loading to the blade and an impulsive noise radiation called blade vortex interaction (BVI) noise. The characteristic signatures of this type of noise differ, depending on whether the interaction occurs on the advancing or retreating sides of the rotor plane. In forward flight, the flow at the tip of the advancing rotor can become locally supersonic. The spatial signature of the shock that closes the supersonic region can radiate to the acoustic far field. This is referred to as high speed impulsive (HSI) noise. See HELICOPTER; SOUND. Philip J. Morris Bibliography. H. H. Hubbard, Aeroacoustics of Flight Vehicles, vols. 1 and 2, Acoustical Society of America, 1995.
Aerodynamic wave drag The force retarding an airplane, especially in supersonic flight, as a consequence of the formation of shock waves. Although the physical laws governing flight at speeds in excess of the speed of sound are the same as those for subsonic flight, the nature of the flow about an airplane and, as a consequence, the various aerodynamic forces and moments acting on the vehicle at these higher speeds differ substantially from those at subsonic speeds. Basically, these variations result from the fact that at supersonic speeds the airplane moves faster than the disturbances of the air produced by the passage of the airplane. These disturbances are propagated at roughly the speed of sound and, as a result, primarily influence only a region behind the vehicle. Causes of wave drag. The primary effect of the change in the nature of the flow at supersonic speeds is a marked increase in the drag, resulting from the formation of shock waves about the configuration. These strong disturbances, which may extend for many miles from the airplane, cause significant
Aerodynamic wave drag
µ
Fig. 1. Shock waves about an airplane at supersonic speeds. µ = angle of obliqueness.
energy losses in the air, the energy being drawn from the airplane. At supersonic flight speeds these waves are swept back obliquely, the angle of obliqueness decreasing with speed (Fig. 1). For the major parts of the shock waves from a well-designed airplane, the angle of obliqueness is equal to sin−1 (1/M), where M is the Mach number, the ratio of the flight velocity to the speed of sound. See SHOCK WAVE; SUPERSONIC FLIGHT. The shock waves are associated with outward diversions of the airflow by the various elements of the airplane. This diversion is caused by the leading and trailing edges of the wing and control surfaces, the nose and aft end of the fuselage, and other parts of the vehicle. Major proportions of these effects also result from the wing incidence required to provide lift. Magnitude of wave drag. Usually the aerodynamic forces acting on a vehicle, such as drag, are considered as coefficients. The wave-drag coefficient for a vehicle depends on many parameters, including the thickness-to-chord ratios and shapes of the airfoil sections of wings and fins, the planform shapes of such surfaces, the length-to-diameter ratio of the body, and the shape of the engine housing. It also depends on the Mach number. For an unswept surface with a thin airfoil whose shape is designed for wave-drag reduction (as discussed below), the coefficient of wave drag is given approximately by Eq. (1), where t/c is 4 Cd = √ (t/c)2 + α 2 M2 − 1
speeds, the wave drag at the zero lift condition is usually more significant than the drag due to wing incidence. But when the Mach number is increased, the relative magnitude of wave drag at the zero lift condition gradually decreases, and the drag associated with wing incidence progressively increases, so at the higher supersonic speeds wave drag due to lift is usually more important than the zero lift value. For a well-designed vehicle, wave drag is usually roughly equal to the sum of the basic skin friction and the induced drag due to lift. See AERODYNAMIC FORCE; AIRFOIL; TRANSONIC FLIGHT. Reduction of wave drag. The wave drag at the zero lift condition is reduced primarily by decreasing the thickness-chord ratios for the wings and control surfaces and by increasing the length-diameter ratios for the fuselage and bodies. Also, the leading edge of the wing and the nose of the fuselage are made relatively sharp (Fig. 2). With such changes, the severity of the diversions of the flow by these elements is reduced, with a resulting reduction of the strength of the associated shock waves. Also, the supersonic drag wave can be reduced by shaping the fuselage and arranging the components on the basis of the area rule. See WING; WING STRUCTURE. The wave drag can also be reduced by sweeping the wing panels (Fig. 3a). Some wings intended for supersonic flight have large amounts of leading-edge sweep and little or no trailing-edge sweep (Fig. 3b). Such planforms are referred to as delta or modified delta. For a simple infinite-span, constant-chord airfoil, the effective Mach number determining the aerodynamic characteristics is the component of the flight Mach number normal to the swept elements (Fig. 4). This Mach number is defined by
subsonic
supersonic
Fig. 2. Comparison of airfoil sections for subsonic flight and supersonic flight.
(1)
the airfoil thickness-to-chord ratio and α is the angle of attack in radians. For surfaces with more complex airfoil shapes and planforms, the computation of wave drag is more complex. For complete vehicles, the wave drag at zero lift is given by the area rule. For supersonic flight, the cross-sectional areas used are obtained in planes inclined at the Mach angle. At the lower supersonic
(a)
(b)
Fig. 3. Wing panels. (a) Sweep-back. (b) Delta.
175
176
Aerodynamics Λ
M flight
M normal = M effective
as that used on subsonic and transonic airplanes, a strong normal shock wave forms ahead of the forward face at supersonic speeds. This shock causes severe loss of energy in the air reaching the engine, and consequent losses of engine performance. In addition, the drag of the airplane is increased. To reduce these losses, special inlets and diffusers which decelerate the airflow to the engine by a series of weak disturbances are used. See SUPERSONIC Richard T. Whitcomb DIFFUSER. Bibliography. J. D. Anderson, Jr., Fundamentals of Aerodynamics, 4th ed., 2005; J. J. Bertin, Aerodynamics for Engineers, 4th ed., 2002; A. M. Kuethe and C.-Y. Chow, Foundations of Aerodynamics: Bases of Aerodynamic Design, 5th ed., 1997.
airfoil
Aerodynamics
Fig. 4. Geometry that determines the effective Mach number for an infinite-span swept airfoil.
Eq. (2), where is the angle of the sweep. If the Mnormal = Mflight × cos
(2)
sweep is such that Mnormal is less than the subsonic Mach number at which the initial onset of wave drag occurs for the unswept airfoil, the swept airfoil will have no wave drag for the flight Mach number. For flight Mach numbers for which the normal component is greater than that for the onset of a shock on the unswept airfoil, the wave drag is not eliminated but is significantly reduced. For practical finite-span swept or delta wings, the reductions of wave drag are usually less than for the ideal infinite-span case. These reductions can be substantially improved by the proper chordwise and spanwise variations of thickness, camber, and incidence. The shape changes required are now determined using very complex fluid-dynamic relationships and supercomputers. See COMPUTATIONAL FLUID DYNAMICS. When the speed is increased to supersonic values, an airplane at a given attitude and altitude experiences large increases in drag, in addition to those associated with the different nature of the flow, because of the higher dynamic pressure at these higher speeds. To offset this effect, supersonic airplanes usually fly at considerably higher altitudes than subsonic vehicles. For example, for efficient flight at Mach 2, an airplane must fly at an altitude of about 60,000 ft (18,000 m). A major problem associated with supersonic flight, particularly at the higher supersonic speeds, is that of taking air into the engines. This air must be decelerated from the flight velocity to a relatively low speed at the compressor of the engine, without excessive energy losses. With a simple inlet, such
The applied science that deals with the dynamics of airflow and the resulting interactions between this airflow and solid boundaries. The solid boundaries may be a body immersed in the airflow, or a duct of some shape through which the air is flowing. Although, strictly speaking, aerodynamics is concerned with the flow of air, now the term has been liberally interpreted as dealing with the flow of gases in general. Depending on its practical objectives, aerodynamics can be subdivided into external and internal aerodynamics. External aerodynamics is concerned with the forces and moments on, and heat transfer to, bodies moving through a fluid (usually air). Examples are the generation of lift, drag, and moments on airfoils, wings, fuselages, engine nacelles, and whole airplane configurations; wind forces on buildings; the lift and drag on automobiles; and the aerodynamic heating of high-speed aerospace vehicles such as the space shuttle. Internal aerodynamics involves the study of flows moving internally through ducts. Examples are the flow properties inside wind tunnels, jet engines, rocket engines, and pipes. In short, aerodynamics is concerned with the detailed physical properties of a flow field and also with the net effect of these properties in generating an aerodynamic force on a body immersed in the flow, as well as heat transfer to the body. See AERODYNAMIC FORCE; AEROTHERMODYNAMICS. Types of aerodynamic flow. Aerodynamics can also be subdivided into various categories depending on the dominant physical aspects of a given flow. In lowdensity flow the characteristic size of the flow field, or a body immersed in the flow, is of the order of a molecular mean free path (the average distance that a molecule moves between collisions with neighboring molecules); while in continuum flow the characteristic size is much greater than the molecular mean free path. More than 99% of all practical aerodynamic flow problems fall within the continuum category. See RAREFIED GAS FLOW. Continuum flow can be subdivided into viscous flow, which is dominated by the dissipative effects of viscosity (friction), thermal conduction, and mass
Aerodynamics diffusion; and inviscid flow, which is, by definition, a flow in which these dissipative effects are negligible. Both viscous and inviscid flows can be subdivided into incompressible flow, in which the density is constant, and compressible flow, in which the density is a variable. In low-speed gas flow, the density variation is small and can be ignored. In contrast, in a high-speed flow the density variation is keyed to temperature and pressure variations, which can be large, so the flow must be treated as compressible. See COMPRESSIBLE FLOW; FLUID FLOW; INCOMPRESSIBLE FLOW; VISCOSITY. Compressible flow regimes. In turn, compressible flow is subdivided into four speed regimes: subsonic flow, transonic flow, supersonic flow, and hypersonic flow. These regimes are distinguished by the value of the Mach number, which is the ratio of the local flow velocity to the local speed of sound (Fig. 1). See MACH NUMBER. Subsonic flow. A flow is subsonic if the Mach number is less than 1 at every point. Subsonic flows are characterized by smooth streamlines with no discontinuity in slope (Fig. 1a). Disturbances (caused by, say, the sudden deflection of the trailing edge of the airfoil) propagate both upstream and downstream and are felt throughout the entire flow field. The flow over light, general-aviation airplanes is subsonic. See SUBSONIC FLIGHT. Transonic flow. A transonic flow is a mixed region of locally subsonic and supersonic flow. The flow far upstream of the airfoil can be subsonic (Fig. 1b), but as the flow moves around the airfoil surface it speeds up, and there can be pockets of locally supersonic flow over both the top and bottom surfaces of the airfoil. If the flow far upstream of the airfoil is at a low supersonic value, say with a Mach number of 1.1 (Fig. 1c), a bow shock wave is formed in front of the body. Behind the nearly normal portion of the shock wave, the flow is locally subsonic. This subsonic flow subsequently expands to a low supersonic value over the airfoil. See SHOCK WAVE; TRANSONIC FLIGHT. Supersonic flow. In a supersonic flow, the local Mach number is greater than 1 everywhere in the flow. Supersonic flows are frequently characterized by the presence of shock waves. Across shock waves, the flow properties and the directions of streamlines change discontinuously (Fig. 1d ), in contrast to the smooth, continuous variations in subsonic flow. If a disturbance is introduced into a steady supersonic flow at some point, the effects of this disturbance are not felt upstream. Instead, these disturbances pile up along a front just in front of the nose, creating a shock wave in the flow. See SUPERSONIC FLIGHT. Hypersonic flow. This is a regime of very high supersonic speeds. A conventional rule is that any flow with a Mach number equal to or greater than 5 is hypersonic (Fig. 1e). Examples include the space shuttle during ascent and reentry into the atmosphere, and the flight of the X-15 experimental vehicle. The kinetic energy of many hypersonic flows is so high that, in regions where the flow velocity decreases, kinetic energy is traded for internal energy of the gas, creating high temperatures. Aerodynamic
M ∞ < 0.8 (a)
M>1 0.8 < M ∞ < 1
M>1
(b) bow shock trailing-edge shock
M1
1< M ∞< 1.2
(c)
shock wave
M>1
expansion wave
M ∞ > 1.2
(d)
M∞ > 5
shock layer
(e) Fig. 1. Regimes for aerodynamic flows; M∞ = Mach number for flow far upstream of the airfoil: (a) subsonic flow; (b) transonic flow with M∞ < 1; (c) transonic flow with M∞ > 1; (d) supersonic flow; and (e) hypersonic flow, where the thin, hot shock layer has viscous interaction and chemical reactions.
heating is a particularly severe problem for bodies immersed in a hypersonic flow. See ATMOSPHERIC ENTRY; HYPERSONIC FLIGHT. Nonaeronautical applications. Applications of aerodynamics cut across various disciplines. For example, automobile aerodynamics has taken on new importance in light of the incentive for energyefficient vehicles. External aerodynamics is important in the reduction of aerodynamic drag on an automobile. With wind-tunnel testing (Fig. 2) and computational aerodynamic analyses, this drag has gradually been reduced by up to 70%, and similar improvements have been reported for trucks and trailers. Internal aerodynamics is important in understanding flows through engine and exhaust systems, thereby contributing to the
177
178
Aeroelasticity
Fig. 2. Aerodynamic flow over an automobile. (Ford Motor Co.)
design of engines with improved gas mileage. In another field, environmental aerodynamics focuses on the study of airflows in heating and airconditioning systems, and on the circulation patterns of air inside buildings and homes. Studies of the external airflow around buildings, especially under hurricane conditions, are of great interest. See COMPUTATIONAL FLUID DYNAMICS; GAS DYNAMICS; WIND TUNNEL. John D. Anderson, Jr. Bibliography. J. D. Anderson, Jr., Fundamentals of Aerodynamics, 4th ed., 2005; J. D. Anderson, Jr., Introduction to Flight, 5th ed., 2005; J. J. Bertin, Aerodynamics for Engineers, 4th ed., 2002; A. M. Kuethe and C. Y. Chow, Foundations of Aerodynamics, 5th ed., 1997.
Aeroelasticity The branch of applied mechanics which deals with the interaction of aerodynamic, inertial, and structural forces. It is important in the design of airplanes, helicopters, missiles, suspension bridges, power lines, tall chimneys, and even stop signs. Variations on the term aeroelasticity have been coined to denote additional significant interactions. Aerothermoelasticity is concerned with effects of aerodynamic heating on aeroelastic behavior in high-speed flight. Aeroservoelasticity deals with the interaction of automatic controls and aeroelastic response and stability. The primary concerns of aeroelasticity include flying qualities (that is, stability and control), flutter, and structural loads arising from maneuvers and atmospheric turbulence. Methods of aeroelastic analysis differ according to the time dependence of the inertial and aerodynamic forces that are involved. For the analysis of flying qualities and maneuvering loads wherein the aerodynamic loads vary relatively slowly, quasi-static methods are applicable, although autopilot interaction could require more general methods. The remaining problems are dynamic, and methods of analysis differ according to whether the time dependence is arbitrary (that is, transient or random) or simply oscillatory in the steady state.
Lift distribution and divergence. The redistribution of airloads caused by structural deformation will change the lifting effectiveness on the aerodynamic surfaces from that of a rigid vehicle. The simultaneous analysis of equilibrium and compatibility among the external airloads, the internal structural and inertial loads, and the total flow disturbance, including the disturbance resulting from structural deformation, leads to a determination of the equilibrium aeroelastic state. If the airloads tend to increase the total flow disturbance, the lift effectiveness is increased; if the airloads decrease the total flow disturbance, the effectiveness decreases. In a limiting case of increasing lift effectiveness, there is a critical speed at which the rate of change of airload with deformation equals the rate of change of structural reaction, and no statically stable equilibrium condition exists; at a higher speed the deformation will increase to a point of structural failure. This critical speed is called the divergence speed. The effectiveness of a lifting surface influences the selection of planform size. Optimum design for strength requires design for the actual loads applied, taking due account of the redistribution of loading resulting from flexibility. The extent of changes in effectiveness and load distribution on a straight wing depends on its torsional stiffness and the chordwise distance between the aerodynamic center of the airfoil and its center of twist. Divergence can occur on a straight wing if the airfoil aerodynamic center is forward of the center of twist. Sweep has a significant influence on effectiveness and redistribution of loads to an extent also determined by the surface bending stiffness, while the effectiveness of stabilizers is also influenced by the stiffness of the fuselage. The bending of a swept-back wing has a substantial stabilizing effect on divergence to the extent that only a slight amount of sweep is required to eliminate the possibility of divergence. The redistribution of wing loading has a slight effect on the induced drag, but if the wing is swept it may have a more significant effect on static longitudinal stability. If the wing is swept-back, the redistribution of loading usually shifts the spanwise center of pressure location inboard, and correspondingly the chordwise center of pressure location moves forward. This forward movement of center of pressure reduces the static stability. The same effect occurs on swept-back stabilizing surfaces and reduces effective tail lengths. This can be an extremely critical problem on a vertical stabilizer since it provides the primary source of static directional stability, whereas it is the combination of wing and horizontal stabilizer that determines the static longitudinal stability. The bending of a wing will also introduce a change in dihedral that may affect the dynamic lateral-directional stability of the aircraft. See WING STRUCTURE. On each of the attempts to attain crewed, powered flight in October and December of 1903, a wing of Samuel Langley’s aerodome collapsed, and the failures may have been caused by torsional divergence. The first successful flight of a power-driven heavierthan-air machine was achieved the week following
Aeroelasticity Langley’s second failure, when the Wright brothers’ biplane flew at Kitty Hawk, North Carolina, on December 17. The Wright biplane experienced no aeroelastic difficulties and, in fact, made use of an aeroelastic control system in which warping of the wing gave rise to the necessary controlling aerodynamic forces. The Wright brothers had also noted a reduction in thrust due to twisting in their experimental work on propeller blades, although it had no significant effect on the blades’ performance. During World War I the Fokker D-8 had just been put into service when three failures occurred in high-speed maneuvers. As a result of static strength and deflection measurements, it became apparent that the failures were due to torsional divergence; the load resulting from the air pressure in a steep dive would increase faster at the wing tips than at the middle. Control effectiveness and reversal. The airloads induced by means of a control-surface deflection also induce an aeroelastic loading of the entire system. Equilibrium is determined as above in the analysis of load redistribution. Again, the effectiveness will differ from that of a rigid system, and may increase or decrease depending on the relationship between the net external loading and the deformation. In a situation of decreasing control-surface effectiveness, the speed at which the effectiveness vanishes is called the reversal speed. At this speed the control-surface– induced air loads are exactly offset by the resulting aeroelastic loads. At higher speeds the aeroelastic loads will exceed the control-surface loading, and the resultant load will act in the reverse direction from the control-surface loading, causing the control system to act in a direction opposite to that desired. Control effectiveness determines the selection of control-surface size and maximum travel required. It may determine the need for all-movable horizontal or vertical stabilizers. Loss in control-surface effectiveness on a straight planform depends on the torsional stiffness of the forward primary surface, to some extent on the torsional stiffness of the control surface as well as its actuation system, and on the distance between the airfoil aerodynamic center and the aerodynamic center of the loading induced by the control surface. Sweepback, which is favorable in the case of divergence, has the unfavorable effect of decreasing control-surface effectiveness. The effectiveness of a tail or canard control surface is also influenced by fuselage bending and, in the case of a vertical stabilizer, torsion of the fuselage. Control reversal is an extremely critical problem with ailerons, and may involve severe weight penalties in providing sufficient torsional stiffness and, in the case of a sweptback wing, bending stiffness. The problem is less critical with elevators and rudders and, of course, is eliminated by the use of all-movable stabilizers for control. Unlike divergence, the phenomenon of reversal is not destructive in itself, but it can lead to structural failure when a marginally stable vehicle cannot be controlled. See AILERON; AIRCRAFT RUDDER; ELEVATOR (AIRCRAFT). Aileron reversal first became a crucial problem for World War II fighters. High roll rates at high speeds
were essential in combat maneuvers, and any loss in rolling ability put a fighter pilot at a serious disadvantage. The designers of the Spitfire, the FW190, and the P-51 sought very high roll rates, and this led to shorter wingspans, aerodynamically balanced ailerons for larger deflection at high speeds, and higher torsional stiffness to increase the reversal speed well beyond the performance limits of the aircraft. The Japanese Zero, on the other hand, was designed for lightness and minimum-radius turns which resulted in a low reversal speed. The U.S. Navy pilots soon discovered this weakness and avoided circling combat but established high-speed, single-pass techniques, where their superior dive speeds and high roll rates could not be followed by the Zero operating close to its aileron reversal speed. Flutter. A self-excited vibration is possible if a disturbance to an aeroelastic system gives rise to unsteady aerodynamic loads such that the ensuing motion can be sustained. At the flutter speed a critical phasing between the motion and the loading permits extraction of an amount of energy from the airstream equal to that dissipated by internal damping during each cycle and thereby sustains a neutrally stable periodic motion. At lower speeds any disturbance will be damped, while at higher speeds, or at least in a range of higher speeds, disturbances will be amplified. The simplest type of flutter occurs when a single motion induces an aerodynamic force having a component in the direction of the motion and in phase with the velocity. This is described as a case of negative aerodynamic damping or single-degreeof-freedom flutter. An example is a power line that “gallops” as a result of ice formation which gives the cross section of the cable something of an airfoil shape. The term classical flutter is used to denote the more complicated instability that typically arises from a critical coupling of two or more modes of motion, each of which is stable by itself. The mechanism for classical flutter may be described in terms of the bending deflection, the angle of twist, and the disturbance in lift. If the structural characteristics (that is, stiffness and inertia) are such that an upward incremental lift, corresponding to a leading-edge upward twist, occurs when the airfoil is bending upward, and, when the motion reverses, the incremental lift acts downward as the airfoil is bending downward, energy will be extracted from the airstream. If the energy extracted in each cycle exceeds that dissipated by internal structural damping, the motion will oscillate with an increasing amplitude until a structural failure occurs. The parameters of the flutter problem are those of the divergence problem with the addition of the dynamic inertial characteristics. The flutter speed is basically proportional to the fundamental torsion frequency so that torsional stiffness is a primary design variable. A forward center of gravity is stabilizing because it tends to drive the bending and twisting motions out of phase, which results in the incremental lift acting in the opposite direction from the bending motion. Mass balance is therefore
179
180
Aeromonas another primary design variable. Flutter speed is usually reduced at high altitudes because aerodynamic damping decreases faster than aerodynamic stiffness (that is, lift effectiveness) with air density. Flutter is essentially a resonance phenomenon between two motions, so that another design consideration is separation of natural vibration frequencies. The effect of sweepback is not as significant on flutter characteristics as it is with divergence, although increasing the sweepback angle does increase the flutter speed. An interesting design possibility is to suppress flutter by automatic control systems, and development of such systems has been undertaken. In order to control the aerodynamic forces, either an additional control surface or a rapidly responding actuator for an existing control surface is required, and an example configuration is a primary surface with both leading- and trailing-edge control surfaces. For practical applications the automatic control system will need multiple redundancies since a component failure could result in violent flutter. The first known case of aerodynamic flutter occurred during World War I on the horizontal tail of a Handley-Page 400, twin-engine bomber. The flutter apparently was caused by coupling of the fuselage torsion mode and the antisymmetrical rotations of the independently actuated right and left elevators. The coupling was eliminated by connecting the elevators to a common actuating torque tube. Toward the end of the war, tail flutter difficulties were overcome by the same redesign, and the attachment of elevators to the same torque tube became a standard feature of subsequent designs. Soon after the war, static mass balancing of ailerons about their hinge lines was found to be an effective means of avoiding wing-aileron flutter. This technique has proved to be both a design and redesign procedure of permanent value, for, ever since, there have been occasional instances of coupled control-surface flutter involving both wings and tail surfaces that did not lead to destruction but could be eliminated merely by increasing the static mass balance of the control surface. See FLUTTER (AERONAUTICS). Gust response. Transient meteorological conditions such as wind shears, vertical drafts, mountain waves, and clear air or storm turbulence impose significant dynamic loads on aircraft. So does buffeting during flight at high angles of attack or at transonic speeds. The response of the aircraft determines the stresses in the structure and the comfort of the occupants. Aeroelastic behavior makes a condition of dynamic overstress possible; in many instances, the amplified stresses can be substantially higher than those that would occur if the structure were much stiffer. See CLEAR-AIR TURBULENCE; LOADS, DYNAMIC; TRANSONIC FLIGHT. Structural design considerations for gusts include both strength and fatigue life. Atmospheric turbulence and buffeting are random phenomena and can be described only statistically, but wind shears and mountain waves take on the appearance of discrete gusts. Since statistical data on all atmospheric conditions are limited, recourse to arbitrary design criteria is necessary. Strength has generally been determined
for response to a discrete gust with a specified profile and maximum velocity normal to the flight path, while fatigue life is estimated by statistical methods based on the power spectral characteristics of the atmosphere. However, strength is frequently also calculated by the statistical methods since static strength may be regarded as the limiting case of fatigue under a single loading. The gust loading condition is a critical one in the design of the wings of large aircraft, and the amount of sweep and the location of wing-mounted engines influence the dynamic response. The basic parameters that determine the response are the longitudinal dynamic flight characteristics of the aircraft, primarily its short-period frequency and damping, and the bending stiffness and frequency of the wing. These same factors, coupled with fuselage flexibility, also determine the aircraft ride qualities. Both gust response loads and ride qualities have been improved by use of automatic control systems. The gustiness of the atmosphere has obviously been a concern throughout the history of flight. The first report of the National Advisory Committee for Aeronautics [now the National Aeronautics and Space Administration (NASA)], published in 1915, was an experimental and theoretical investigation of aircraft response to gusts. The first practical applications of automatic controls for gust alleviation and structural dynamic stability augmentation were made to the B-52 and B-70 bombers. Successful demonstrations of the load alleviation systems were begun with B-52 flights in late 1967 and B-70 flights in 1968. See AERODYNAMICS; FLIGHT CHARACTERISTICS; SUBSONIC FLIGHT. William P. Rodden Bibliography. R. L. Bisplinghoff and H. Ashley, Principles of Aeroelasticity, 2d ed., 1962, reprint 2002; E. H. Dowell (ed.), A Modern Course in Aeroelasticity, 4th ed., 2004; Y.-C. B. Fung, An Introduction to the Theory of Aeroelasticity, 1955, reprint 2002; W. P. Jones (ed.), AGARD Manual on Aeroelasticity, 1960–1968; D. McRuer, I. Ashkenas, and D. Graham, Aircraft Dynamics and Automatic Control, 1973.
Aeromonas A bacterial genus in the family Vibrionaceae comprising oxidase-positive, facultatively anaerobic, monotrichously flagellated gram-negative rods. The mesophilic species are A. hydrophila, A. caviae and A. sobria; the psychrophilic one is A. salmonicida. Aeromonads are of aquatic origin and are found in surface and wastewater but not in seawater. They infect chiefly cold-blooded animals such as fishes, reptiles, and amphibians and only occasionally warmblooded animals and humans. Human wound infections may occur following contact with contaminated water. Septicemia has been observed mostly in patients with abnormally low white blood counts or liver disease. There is evidence of intestinal carriers. The three mesophilic species are also associated with diarrheal disease (enteritis and colitis) worldwide. A heat-labile
Aeronautical engineering enterotoxin and cytotoxin(s) are produced only by A. hydrophila and A. sobria, but experimental diarrhea could not be produced in humans. See DIARRHEA. A related lophotrichous genus, Plesiomonas (single species, P. shigelloides), is also known as an aquatic bacterium and is associated with diarrhea chiefly in subtropical and tropical areas. It is also found in many warm-blooded animals. Systemic disease in humans is rare. See MEDICAL BACTERIOLOGY. Alexander von Graevenitz Bibliography. E. H. Lennette et al. (eds.), Manual of Clinical Microbiology, 4th ed., 1985.
Aeronautical engineering That branch of engineering concerned primarily with the special problems of flight and other modes of transportation involving a heavy reliance on aerodynamics or fluid mechanics. The main emphasis is on airplane and missile flight, but aeronautical engineers work in many related fields such as hydrofoils, which have many problems in common with aircraft wings, and with such devices as air-cushion vehicles, which make use of airflow around the base to lift the vehicle a few feet off the ground, whereupon it is propelled forward by use of propellers or gas turbines. See AERODYNAMICS; AIRPLANE; FLUID MECHANICS; HYDROFOIL CRAFT. Aeronautical engineering expanded dramatically after 1940. Flight speeds increased from a few hundred miles per hour to satellite and space-vehicle velocities. The common means of propulsion changed from propellers to turboprops, turbojets, ramjets, and rockets. This change gave rise to new applications of basic science to the field and a higher reliance on theory and high-speed computers in design and testing, since it was often not feasible to proceed by experimental methods only. See JET PROPULSION; PROPULSION; ROCKET PROPULSION; SPACE TECHNOLOGY. Aeronautical engineers frequently serve as system integrators of important parts of a design. For example, the control system of an aircraft involves, among other considerations, aerodynamic input from flow calculations and wind-tunnel tests; the structural design of the aircraft (since the flexibility and strength of the structure must be allowed for); the mechanical design of the control system itself; electrical components, such as servo-mechanisms; hydraulic components, such as hydraulic boosters; and interactions with other systems that affect the control of the aircraft, such as the propulsion system. The aeronautical engineer is responsible for ensuring that all of these factors operate smoothly together. The advent of long-range ballistic missiles and spacecraft presented new challenges for aeronautical engineering. Flight velocities increased up to 35,000 ft/s (10,668 m/s), and the resulting heating problems required new kinds of materials and cooling-system concepts to protect the spacecraft or reentry vehicle. For example, materials were developed by ground testing in high-temperature electric
arcs and rocket exhausts and then were tested by flights of crewless models at full speed. The resulting heat shields were successfully used to return astronauts from the Moon and to protect the warheads of ballistic missiles while maintaining missile accuracy. See BALLISTIC MISSILE. The techniques that aeronautical engineers had developed for aircraft were applied to help develop such devices as the air-cushion vehicle, which is of much interest to the military for landing troops on a beach. Helicopters and vertical takeoff and landing aircraft are other developments in which aeronautical engineering has played an important role. Hydrofoils are basically boats that fly on wings under the water. Because the water is so dense, these wings are usually small compared to aircraft wings, but the same principles apply. The control of such wings is very sensitive, and careful attention must be paid to the control system used lest the high forces involved cause severe damage to the wings or to the boat itself. See AIR-CUSHION VEHICLE; HELICOPTER; VERTICAL TAKEOFF AND LANDING (VTOL). Aerodynamics has always been one of the main tools of aeronautical engineering. As noted above, aerodynamics interacts closely with control design. For example, as an aircraft approaches the ground to land, the effect of the solid barrier of the ground affects the aerodynamics of the aircraft. This effect can be simulated in a computer model by placing another wing in a mirror image under the ground. The resulting flow pattern corresponds to the real wing and the solid ground. In any but the simplest cases, this type of calculation cannot be carried out except with the aid of high-speed computers, which have made possible the calculation of entire flow fields around aircraft and spacecraft as they pass through the atmosphere. Wind-tunnel testing of scale models is another important source of aerodynamic data. See COMPUTATIONAL FLUID DYNAMICS; WIND TUNNEL. Aircraft and missile structural engineers have raised the technique of designing complex structures to a level never considered possible before the advent of high-speed computers. Structures can now be analyzed in great detail and the results incorporated directly into computer-aided design (CAD) programs. These programs are now so complete that in some cases no drafting work with paper and pencil is used at all and all drawings are printed directly from the design devised by the engineer working only on the computer screen. See COMPUTER-AIDED DESIGN AND MANUFACTURING. These examples illustrate the scope of aeronautical engineering. In general, it overlaps many other disciplines, and the aeronautical engineer may concentrate on aerodynamics, structures, control systems, propulsion systems, or the overall design of the vehicle. John R. Sellars Bibliography. J. E. Allen, Aerodynamics, 1982; W. Biddle, Barons of the Sky: From Early Flight to Strategic Warfare, 1991; P. A. Hanle, Bringing Aerodynamics to America: A History of Aerodynamics Research, 1982; F. E. Weick, From the Ground Up: The Autobiography of an Aeronautical Engineer, 1988.
181
1022 1024 1026
1024
1024
L BIS
DLH U
u
The branch of meteorology that deals with atmospheric effects on the operation of vehicles in the atmosphere, including winged aircraft, lighter-thanair devices such as dirigibles, rockets, missiles, and projectiles. The air which supports flight or is traversed on the way to outer space contains many potential hazards. Nine major hazards are considered in this article, which is principally concerned with commercial and general aviation. Low visibility. Poor visibility caused by fog, snow, dust, and rain is a major cause of aircraft accidents and the principal cause of flight cancellations or delays. For the average private pilot with a minimum of equipment in the aircraft, an unexpected encounter with low-visibility conditions can be a serious matter, although technology is bringing aviation closer to the hitherto elusive goal of all-weather flying. The weather conditions of ceiling and visibility required by regulations for crewed aircraft during landing or takeoff are determined by electronic and visual aids operated by the airport and installed in the aircraft. There are three categories of airports as related to the installation of various levels of specialized aids to navigation. At category 1 airports, landings and takeoffs are permitted when the ceiling is 200 ft (60 m) or more and the runway visual range (RVR) is 1800 ft (550 m) or more. The pilot operates the controls at touchdown, although the initial approach may be controlled by instruments. At category 2 airports, landings and takeoffs are authorized down to a ceiling height of about 150 ft (45 m) and runway visual range of about 500 ft (150 m). At category 3 airports, sophisticated aids to navigation, such as Global Positioning System (GPS), provide for safe, fully automated, and electronically coupled landings under “zero, zero” conditions. Airports smaller than the major city airports may have higher ceiling requirements, as may airports with tall buildings or mountains nearby. See AIR NAVIGATION; AIRCRAFT INSTRUMENTATION; INSTRUMENT LANDING SYSTEM (ILS). The accurate forecasting of terminal conditions is critical to flight economy, and to safety where sophisticated landing aids are not available. Improved prediction methods are under continuing investigation and development, and are based on mesoscale and microscale meteorological analyses, electronic computer calculations, radar observations of precipitation areas, and observations of fog trends. See MESOMETEOROLOGY; MICROMETEOROLOGY; RADAR METEOROLOGY. Turbulence and low-altitude wind shear. Atmospheric turbulence is principally represented in vertical currents and their departures from steady, horizontal airflow. When encountered by an aircraft, turbulence produces abrupt excursions in aircraft position, sometimes resulting in discomfort or injury to passengers, and sometimes even structural damage or failure. Major origins of turbulence are (1) mechanical, caused by irregular terrain below the flow of air; (2) thermal, associated with vertical
WEATHER MAP TYPES FOR APRIL MAP FOR 0630Z APRIL 20 H 1026
12
Aeronautical meteorology
1012 1014 1016 1018 1020
Aeronautical meteorology
10
182
1014
CHI
MLI
OMA U
CLE
22
10
DCA SDF
ICT
H
SYR LGA
ROA
ORF
L
1016
MEM
AMA
SHV
1016
10 14 10 12
L
ATL
1020
U
U 1012
H
MSY
1014
1018
1016
Fig. 1. Severe warm-front thunderstorm situation over the eastern United States. Open or solid circles with three letters represent reporting stations, such as OMA for Omaha and CHI for Chicago. Contour lines express millibar pressure. (United Air Lines)
currents produced by heating of air in contact with the Earth’s surface; (3) thunderstorms and other convective clouds (Fig. 1); (4) mountain wave, a regime of disturbed airflow leeward of mountains or hills, often comprising both smooth and breaking waves formed when stable air is forced to ascend over the mountains; and (5) wind shear, usually variations of horizontal wind in the vertical direction, occurring along air-mass boundaries, temperature inversions (including the tropopause), and in and near the jet stream. Vertical and horizontal wind gusts in the upper troposphere and lower stratosphere, as encountered by jet transport aircraft that frequent these altitudes, are called clear-air turbulence. Turbulence encounters are the principal cause of nonground impact incidences in commercial aviation. The distinction between clear air and thunderstorm turbulence becomes uncertain as aircraft are routinely operated at altitudes above 30,000 ft (9 km) and encounter thunderstorms growing to more than 60,000 ft (18 km). While encounters with strong turbulence anywhere in the atmosphere represent substantial inconvenience, encounters with rapid changes in wind speed and direction at low altitude can be catastrophic. Generally, wind shear is most dangerous when encountered below 1000 ft (300 m) above the ground, where it is identified as low-altitude wind shear. Intense convective microbursts, downdrafts usually associated with thunderstorms (Fig. 2), have caused many aircraft accidents often resulting in a great loss of life. The downdraft emanating from convective clouds, when nearing the Earth’s surface, spreads horizontally as outrushing rain-cooled air. When entering a microburst outflow, an aircraft first meets a headwind that produces increased
Aeronautical meteorology
2.4 (4) dry air inflow
1.8 (3) microburst downdraft
aircraft approach path
warm air inflow
headwind
12 (20)
1.2 (2) headwind
9 (15)
tailwind 6 (10) 3 (5) distance, mi (km)
gust front
0.6 (1)
height, mi (km)
performance by way of increased airspeed over the wings. Then within about 5 s, the aircraft encounters a downdraft and then a tailwind with decreased performance. A large proportion of microburst accidents, both after takeoff and on approach to landing, are caused by this performance decrease, which can result in rapid descent. See THUNDERSTORM. Turbulence and low-altitude wind shear can readily be detected by a special type of weather radar, termed Doppler radar. By measuring the phase shift of radiation backscattered by hydrometeors and other targets in the atmosphere, both turbulence and wind shear can be clearly identified. It is anticipated that Doppler radars located at airports, combined with more thorough pilot training regarding the need to avoid microburst wind shear, will provide desired protection from this dangerous aviation weather phenomenon. See DOPPLER RADAR; METEOROLOGICAL RADAR. Hail. Fabric-covered aircraft are particularly susceptible to hail damage. The advent of all-metal aircraft reduced the probability of damage from small (0.2 in. or 5 mm diameter) hail, but larger hail (1–2 in. or 25–50 mm diameter) can be very destructive—shattering the wind screen, dimpling skin surfaces, and cracking radomes. Avoidance of thunderstorms is the best preventive when in flight; on-board radar is a valuable detection device, especially at night and for storms embedded in widespread clouds. See HAIL. Turbulence. Turbulence encounters in and near thunderstorms are increasing with attendant increasing reports of injury to passengers. Part of the increase is because more aircraft fly above 30,000 ft (9 km) and sky space is more crowded. Thunderstorms extend above 60,000 ft (18 km) with their upper portions and anvils spreading out for tens of miles (or kilometers). These storms also compete for space with air traffic. Turbulence at these levels is often located in the regions between the high reflectivity volumes—that is, on the edges of the larger scale vertical motion cores—and in areas where lightning activity is less frequent (charge separation tends to be destroyed by turbulence). Severe turbulence is encountered as much as 20 miles (37 km) from the thundersorm cores. Turbulence associated with mountains and wind shears is frequently encountered in clear air, with little forewarning to the pilot. The unexpected nature of the clear-air turbulence adds to its danger, especially with respect to passenger safety. The intensity of clear-air turbulence depends on the size and speed of the aircraft. Generally, the accelerations are larger at faster speeds. Clear-air turbulence exists even at the 49,000–65,000 ft (15,000–20,000 m) flight levels of supersonic transport aircraft, and the paths of such aircraft must be planned to avoid areas where strong jet streams are located over tall mountains. Since clear-air turbulence is a small-scale phenomenon (patches of turbulence frequently measure only a few miles in diameter and a few hundred to a few thousand feet in thickness) the upper-air sounding network is insufficient for locating or predicting
183
0
Fig. 2. Schematic of a thunderstorm intersecting the approach path to an airport. Arrows ◦ indicate wind-flow pattern. Broken line is 3 glide slope approach. Heavy lines outline the leading edge of the thunderstorm’s cold air outflow (gust front). Note rapid change in wind in center of microburst downdraft.
the precise locations of clear-air turbulence zones. Remote sensors, such as the Doppler lidar and the infrared radiometer, are able to detect zones of wind shear and temperature inversions ahead of the aircraft, permitting evasive actions to be taken or passenger warning signs to be activated. See CLEAR-AIR TURBULENCE; LIDAR. Upper winds and temperature. Since an aircraft’s speed is given by a propulsive component plus the speed of the air current bearing the aircraft, there are aiding or retarding effects depending on wind direction in relation to the track flown. Wind direction and speed vary only moderately from day to day and from winter to summer in certain parts of the world, but fluctuations of the vector wind at middle and high latitudes in the troposphere and lower stratosphere can exceed 200 knots (100 m · s−1). Careful planning of long-range flights is completed with the aid of prognostic upper-air charts. These charts are prepared through the use of computers to solve equations that model atmospheric motion. The selection of flight tracks and cruising altitudes having relatively favorable winds and temperature results in lower elapsed flight times, in some cases even with considerable added ground distance. The role of the aeronautical meteorologist is to provide accurate forecasts of the wind and temperature field, in space and time, through the operational ranges of each aircraft involved. For civil jet-powered aircraft, the optimum flight plan must always represent a compromise among wind, temperature, and turbulence conditions. See UPPER-ATMOSPHERE DYNAMICS; WIND. Jet stream. This is a meandering, shifting current of relatively swift wind flow which is embedded in the general westerly circulation at upper levels. Sometimes girdling the globe at middle and subtropical latitudes, where the strongest jets are found, this band of strong winds, generally 180–300 mi (300–500 km) in width, has great operational significance for aircraft flying at cruising levels of 4–9 mi (6–15 km). It causes some serious flight-schedule delays, and unexpected encounters may force unscheduled fuel stops. Average speed in the core of a well-developed jet near the tropopause at middle latitudes in winter is
Aeronautical meteorology 17200 17400
16800 17000
L
16800 17000 17000
17000 17400 17200
L
17 17 80 60 0 0
184
17600
RAP
17800
0
BOS 180000
westbound
18000
CHI
_35
PIT DCA
PUB 18200 18400
ICT
TYS LIT
ELP 18600
+125
EAM JET STR 0
1880
00
190
0
20
19
2 18 0 40 8 1 00 6 18 0 80 18 0 0 0 19 0 20 19
eastbound H
Fig. 3. Use of minimal time track in flight planning over the eastern United States. The solid circles represent weather-reporting places, landing places, or both, such as CHI for Chicago and BOS for Boston. The plus or minus figures shown in the large circles indicate the relation of air speed to ground speed (mi/h) to be expected in the direction of flight. Contour lines express height of the 500-millibar pressure surface in feet. (United Air Lines)
close to 100 knots (50 m · s−1), but speeds nearly twice that are fairly common, and extremes may be three times larger. The jet stream challenges the forecaster and the flight planner to utilize tailwinds to the greatest extent possible on downwind flights and to avoid retarding headwinds as much as practicable on upwind flights (Fig. 3). As with considerations of wind and temperature, altitude and horizontal coordinates are considered in flight planning for jetstream conditions. Turbulence in the vicinity of the jet stream is also a forecasting problem. There is also the “low-level” jet, a relatively narrow (100-nautical-mile or 185-km) stream of air having wind speeds of more than 40 knots (20 m · s−1) through a depth of around 1000 ft (300 m) at an altitude of 1000–3000 ft (300–900 m) above the surface. The northward-flowing jet is usually found from Texas on into the Great Lakes Region. When it occurs, the jet brings moist air rapidly northward and greatly influences the formation of convective storms. See JET STREAM. Tropopause. This boundary between troposphere and stratosphere is generally defined in terms of the zone in the vertical air-mass soundings where the temperature lapse rate becomes less than 3◦F/mi (2◦C/km). It is sometimes defined as a discontinuity in the wind shear associated with horizontal temperature gradients. Sometimes the tropopause appears as a layer of irregular temperature variations, caused by overlapping “tropical” and “polar” leaves. The significance of the tropopause to aviation arises from its frequent association with mild forms of clear-air turbulence and change of vertical temperature gradient with altitude. In midlatitudes, the tropopause is commonly located near 6–7 mi (10– 12 km), which is about the cruising altitude of most jet aircraft, while in tropical regions it is around 9– 10 mi (16–17 km), that is, a little below the cruising altitude of supersonic aircraft. The tropopause marks the usual vertical limit of clouds and storms; however, thunderstorms are known to punch through
the surface to heights near 12 mi (20 km) at times, and lenticular clouds showing wave motions strongly developed over mountains have been observed at similar heights. See TROPOPAUSE. Lightning. An electrical discharge to or from an aircraft is experienced as a blinding flash and a muffled explosive sound, usually audible above the roar of the engines. Structural damage to the metal portions of the craft is commonly limited to melting of small spots in the outer skin at the point of entry or exit, and fusing of antenna wires. Nonmetallic surfaces such as radomes and composite tail surfaces may suffer punctures or delaminations. Atmospheric conditions favorable for lightning strikes follow a consistent pattern, characterized by solid clouds or enough clouds for the aircraft to be flying intermittently on instruments; active precipitation of an icy character; and ambient air temperature near or below 32◦F (0◦C). Saint Elmo’s fire, radio static, and choppy air often precede the strike. However, the charge separation processes necessary for the production of strong electrical fields is destroyed by strong turbulence. Thus turbulence and lightning usually do not coexist in the same space. An aircraft usually triggers the lightning stroke, and evasive action by the pilot usually consists of air-speed reduction, change of course as suggested by radar echoes, or change of altitude. Pilots unable to avoid a typical strike situation should turn up cockpit lights and avert their eyes in order to decrease the risk of being blinded. As metal becomes replaced by composite semiconducting materials in aircraft, there is increased risk of damage from lightning, both to the composite section itself and to inadequately shielded microelectric components. See ATMOSPHERIC ELECTRICITY; LIGHTNING; SAINT ELMO’S FIRE. Icing. Modern aircraft operation finds icing to be a major factor in the safe conduction of a flight both in the takeoff and landing phases as well as the in-flight portion of the trip. Ice and frost formation on an aircraft result in increased drag, decreased lift due to disruption of airflow over the wing and tail surfaces, and increased weight. Icing usually occurs when the air temperature is near or below freezing (32◦F or 0◦C) and the relative humidity is 80% or more. Clear ice is most likely to form when the air temperature is between 32 and −4◦F (0 and −20◦C) and the liquid water content of the air is high (large drops or many small drops). As these drops impinge on the skin of an aircraft, the surface temperature of which is 32◦F (0◦C) or less, the water freezes into a hard, high-density solid. When the liquid water content is small and when snow or ice pellets may also be present, the resulting rime ice formation is composed of drops and encapsulated air, producing an ice that is less dense and opaque in appearance. The accumulation of ice can occur in a very short time interval. In flight, the use of deicing boots (inflatable rubber sleeves) or heated air of the leading edge of the wing or tail surfaces can reduce the problem. On the ground, ice and snow can be removed either by mechanical means or by application of
Aeronomy
Aeronomy concentrates on the regions above the tropopause or upper part of the atmosphere. The region of the atmosphere below the tropopause is the site of most of the weather phenomena that so directly affect all life on the planet; this region has primarily been the domain of meteorology. Aeronomy describes the chemical and physical properties of planetary atmospheres and the changes that result from external and internal forces. The results of this interaction impact all altitudes and global distributions of atoms, molecules, ions, and electrons, both in composition and in density. Dynamical effects are seen in vertical and horizontal atmospheric motion, and energy is transferred through radiation, chemistry, conduction, convection, and wave propagation. Thus aeronomy is a truly interdisciplinary field encompassing chemistry, physics, electrical engineering, and mathematics and dealing with all aspects of theory, simulations, modeling, and laboratory and field experiments. Research may involve the use of specialized ground-based optical and radar facilities; satellite, rocket, and balloon instrumentation; and state-of-the-art computational facilities, locally, regionally, and globally. Since aeronomy is not restricted to earth science, optical and radar remote-sensing techniques are also used to study the atmospheres of other planets in the solar system. See METEOROLOGICAL RADAR; REMOTE SENSING. The atmosphere of the Earth is separated into regions defined by the variation of temperature with height (see illustration). Above the surface, temperature minima occur near 10 km (6 mi; the tropopause) and 90 km (54 mi; the mesopause). In between these two minima, the maximum in temperature near 50 km (30 mi) is called the stratopause. The region below the tropopause is the troposphere; the regions between the minimum and maximum
1000 exosphere 1012 1014 1015 thermosphere
100
mesopause mesosphere stratopause
1016 1019 1021 1022 1023
stratosphere
number density/m3
Aeronomy
values are the stratosphere and mesosphere. Above the mesopause the region is called the thermosphere. The temperature gradually increases above the mesopause throughout the thermosphere to a constant value near 300–500 km (180–300 mi). Above this altitude the region is called the exosphere, which extends out to the limit of the detectable terrestrial atmosphere. The region between the tropopause and the mesopause (roughly 10–100 km or 6–60 mi) is called the middle atmosphere, and the region above (>100 km or 60 mi) is called the upper atmosphere. See MESOSPHERE; METEOROLOGY; STRATOSPHERE; TROPOSPHERE. Middle atmosphere. In the middle atmosphere, that region of the atmosphere between the tropopause and the mesopause (10–100 km or 6–60 mi), the temperature varies from 250 K (−9.7◦F) at the tropopause to 300 K (80◦F) at the stratopause and back down to 200 K (−100◦F) at the mesopause (see illus.). These temperatures are average values, and they vary with season and heat and winds due to the effect of the Sun on the atmosphere. Over this same height interval the atmospheric density varies by over five orders of magnitude. Although there is a constant mean molecular weight over this region, that is, a constant relative abundance of the major atmospheric constituents of molecular oxygen and nitrogen, there are a number of minor constituents that have a profound influence on the biosphere and an increasing influence on the change in the atmosphere below the tropopause associated with the general topic of global change. These constituents, called the greenhouse gases (water vapor, ozone, carbon dioxide, methane, chlorine compounds,
altitude, km
deicing solutions. Commercial operations use deicing solutions before takeoff to remove and temporally retard the formation of these increased dragproducing and decreased lift-producing hazards. Carburetor-equipped aircraft may also experience carburetor icing. The sudden cooling of moist air caused by the vaporization of fuel and the expansion of air as it passes through the carburetor can be as much as 60◦F (33◦C) in a few seconds. Water thus condensed is deposited as frost or ice inside the carburetor, thereby restricting the airflow. The result is a power decrease and in some cases even engine failure. Carburetor heating judicially used can reduce this problem. Icy or deeply snow-covered airport runways offer operational problems for aviation, but these factors are reduced in importance through the use of modern snow removal equipment. However, accurate forecasts and accurate delineation of freezing conditions are essential for safe aircraft operations. J. T. Lee; J. McCarthy Bibliography. D. R. MacGorman and W. D. Rust, The Electoral Nature of Storms, 1998.
1024
tropopause troposphere
10
1025 0
1000 500 temperature, K
1500
Typical atmospheric temperature variation with altitude and constituent number density (number of atoms and molecules per cubic meter) at high solar activity. ◦ F = (K × 1.8) − 459.67. 1 km = 0.6 mi.
185
186
Aeronomy nitrogen oxides, chlorofluorocarbons, and others), are all within the middle atmosphere. To fully understand and model this region of the atmosphere requires consideration of the solar variability, horizontal and vertical wind systems, and the reaction rates and time constants involved with all the chemical and radiative processes that control these constituent distributions. This information needs to be obtained simultaneously at all altitudes and all global locations. The magnitude of this task is huge, and thus aeronomy requires national and international coordinated programs. See GREENHOUSE EFFECT. Solar radiation. Solar radiation is the primary energy source that drives all of the various processes in middle-atmospheric aeronomy. Although the solar spectrum can at times be approximated by a body radiating at a temperature around 6000 K (10,300◦F), a detailed study of the spectrum as received at different altitudes in the middle atmosphere shows that absorption by the constituents of the atmosphere modifies the solar radiation at specific wavelengths. Since ozone in the stratosphere strongly absorbs solar radiation below 300 nanometers, the amount of that radiation reaching the ground is very dependent on the total amount of ozone. The harmful relationship between radiation below 300 nm and the Earth’s biosystems is the driving concern about the destruction of ozone through the addition of new chemicals that will modify the historical chemical balance of the stratosphere. See SOLAR RADIATION; STRATOSPHERIC OZONE. Dynamics. Understanding the aeronomy of the middle atmosphere requires the study of the physical motion of the atmosphere. The particular composition and chemistry at any given time, location, or altitude depends on how the various constituents are transported from one region to another. Thus, global circulation models for both horizontal and vertical motions are needed to completely specify the chemical state of the atmosphere. In understanding the dynamics of the middle atmosphere, internal gravity waves and acoustic gravity waves play significant roles at different altitudes, depending on whether the particle motion associated with the wave is purely transverse to the direction of propagation or has some longitudinal component. These waves originate primarily in meteorological events such as wind shears, turbulent storms, and weather fronts; and their magnitude can also depend on orographic features on the Earth’s surface. One example of coupling between regions in the middle atmosphere is seen in the events known as stratospheric warmings, which are produced by gravity waves generated from the large tropospheric waves that warm the stratosphere and cool the mesosphere. In the polar regions the winter circulation tends to be very stable and circumpolar, and it is usually referred to as the polar vortex. In the Antarctic this region is so stable that photolysis of many of the middle atmospheric species does not occur, thus allowing the buildup of species that would otherwise never develop significant concentrations. Upon reexposure to sunlight during the spring,
there is a dramatic photolytic conversion to high concentrations of other species that combine in the presence of polar stratospheric clouds to produce the observed spring ozone depletion (ozone hole) in the lower stratosphere. In the summer polar regions the mesopause temperature becomes very low ( 1
shock wave
body surface
radiation from hot gas
gas boundary layer shock layer Heat transfer in stagnation region of blunted nose cap.
197
198
Affective disorders as a source of electromagnetic noise, which can also adversely affect communication with the vehicle. At still higher temperatures the gas ionizes, producing electromagnetic forces in addition to the radiation. See MAGNETOHYDRODYNAMICS. Aerodynamic heating. Aerodynamic heating is a severe problem at high flight speeds and at all altitudes at which air forms a continuum, that is, up to approximately the altitude at which the mean free path is of the same order of magnitude as the diameter of the vehicle’s nose. Such an altitude might be about 350,000 ft (107 km) for a nose having a diameter of about 1 ft (0.3 m). To investigate aerodynamic heating on a missile, aircraft, spacecraft, or other flying object, one needs detailed information about flight profiles, such as the altitude, velocity, angle of attack, and bank-angle history of the vehicle as a function of time. For most supersonic aircraft, the thermal design can be based on the fixed-equilibrium skin temperature (the temperature at which heat input equals heat output). For a missile, equilibrium is not reached, and the transient nature of missile flight trajectories complicates the temperature analysis. For an atmospheric-entry spacecraft, skin temperature will in some cases reach equilibrium because of its longer flight time, but the equilibrium skin temperature may change with time because of changes in altitude, speed, angle of attack, and bank angle. See ATMOSPHERIC ENTRY. The atmospheric-entry characteristics of spacecraft and of ballistic missiles differ in both the total heat load experienced and in maximum heating rates. Total heat load is greater for the spacecraft, whereas the maximum heating rate is greater for the ballistic missile because of its steeper entry path. For spacecraft with a high lift-to-drag ratio (L /D ≈ 3), the severity of heating can be reduced by trajectory selection, configuration design, and radiation cooling. For spacecraft with medium and low liftto-drag ratios (L /D ≈ 2 or less) and for the ballistic entry vehicle, ablation cooling techniques have been used extensively to protect interior structure. Changing the nose configuration alone cannot completely solve the structural heating problem for a ballistic missile. See SPACECRAFT STRUCTURE. In general, heating analysis is conducted by dividing the vehicle into components; appropriate heat-transfer formulas are then applied to each component to estimate its heating rate. Analysis can be further complicated by boundary-layer transition, boundary-layer interaction, separation flow, shock impingement, mass transfer, nonequilibrium phenomena, gap effects, and rough or wavy surfaces. Accurate estimates of aerodynamic heating also require knowledge of fluid properties of the gas, such as composition, thermal and transport properties, reaction rates, relaxation times, and radiation characteristics within the fluid around the hypersonic vehicle. Under standard conditions air is composed chiefly of diatomic molecules. Motions of translation and rotation are excited at room temperature. As the temperature of air increases, vibration begins, increasing in amplitude as temperature increases,
until eventually the intramolecular bond is broken. The molecule is then said to be dissociated. Still higher energy levels are excited as the temperature increases further; eventually, electrons are separated from their parent atoms and the gas becomes ionized. Shih-Yuan Chen Bibliography. J. D. Anderson, Jr., Hypersonic and High Temperature Gas Dynamics, 1989, reprint, 2000; J. D. Anderson, Jr., Modern Compressible Flow, 3d ed., 2003; J. J. Bertin, Hypersonic Aerothermodynamics, 1994; G. Oates, Aerothermodynamics of Gas Turbine and Rocket Propulsion, 3d ed., 1997; C. Park, Nonequilibrium Hypersonic Aerothermodynamics, 1990.
Affective disorders A group of psychiatric conditions, also known as mood disorders, that are characterized by disturbances of affect, emotion, thinking, and behavior. Depression is the most common of these disorders; about 10–20% of those affected by depression also experience manic episodes (hence this condition is known as manic-depression or bipolar affective disorder). The affective disorders are not distinct diseases but psychiatric syndromes that likely have multiple or complex etiologies. Clinical Syndromes Affective disorders are distinguished from states of normal sadness or elation by several features. Their clinical syndromes consist of characteristic changes in brain and psychological function that affect sleep, appetite, speed of movement, energy, concentration, self-esteem, motivation, and hedonic (pleasure) capacity. During an episode, these features are persistent, occurring day after day for at least several weeks. Affective syndromes also differ from normal mood states by the degree of associated impairment of functioning. The most severe episodes include hallucinations or delusions (unshakable false beliefs of obvious personal significance, typically reflecting themes that are congruent with the mood disturbance). A natural parallel to the clinical syndrome of depression is the state of grief. Grief is associated with many of the symptoms of depression, although the changes of mood and behavior that are associated with grief are less persistent and more circumscribed. Conversely, new romantic love or an infatuation is a natural parallel of a hypomanic episode (that is, a mild episode of mania), although again the normal “affliction” is less persistent and pervasive and should not cause functional impairment. Unresolved grief can progress into clinical depression, however, and there is a pathological form of delusional romantic love that is called erotomania. Aside from the death of a loved one, other types of losses, including divorce, the breakup of a romantic relationship, or setbacks in vocational pursuits, commonly precede the onset of an initial depressive episode. A transitional diagnosis of an adjustment disorder is used when the duration or severity of
Affective disorders depressive symptoms has not crossed the threshold necessary for diagnosis of a depressive disorder. A majority of adjustment disorders, like normal grief, remit as the individual adapts to the stressor. Major depressive disorders. The most common form of affective disorder is a major depressive episode. The episode is defined by a pervasively depressed or low mood (experienced most of the day over a period of 2 weeks or longer) and at least four associated symptoms affecting sleep, appetite, hedonic capacity, interest, and behavior. Major depressive episodes have several clinical forms. Melancholia is a severe episode characterized by anhedonia (inability to experience pleasure), markedly diminished appetite (anorexia; note the anorexia of depression is not related to anorexia nervosa), with weight loss, early morning awakening, observable motor disturbances (extreme slowing, or retardation, or pacing and stereotypic agitated behaviors), and diurnal mood variation (mood is worse in the morning). Some episodes of melancholia are associated with delusions, typically of guilt, sin, illness, punishment, or nihilism. Common among young patients, especially women, is a milder syndrome historically referred to as atypical depression. Atypical depression is characterized by a less pervasive mood disturbance, with mood reactivity preserved (that is, the patient’s spirits go up or down in response to day-to-day events) and by so-called reverse symptoms: oversleeping, overeating, or gaining weight. Significant anxiety symptoms, including phobias and panic attacks, also are common in atypical depression. Both atypical and other milder forms of major depressive disorder were formerly referred to as reactive or neurotic because psychosocial factors often appeared to play a prominent role in their onset. Melancholia, by contrast, is more autonomous and less clearly linked to life stress. In fact, melancholia was formerly called endogenous depression because the cause was presumed to be within the brain. Current thinking on the roles of life events and brainrelated processes is less simplistic: it is likely that both sets of factors are relevant to all forms of affective disorder. Dysthymia. A more chronic, insidious form of depression known as dysthymia “smolders” at a subsyndromal level (that is, there are three or four daily symptoms rather than five or more) for at least 2 years. Dysthymia often begins early in life and, historically, has been intertwined with atypical and neurotic characteristics of depression. There may be a depressive personality or temperament that underlies childhood-onset dysthymia. Most people with untreated dysthymia eventually develop major depressive episodes and, despite its milder symptoms, dysthymia causes considerable impairment. Mania and bipolar affective disorder. About 10–20% of people who experience major depressive episodes will also suffer from abnormal periods of elation, which are referred to as manias (or in less severe form, hypomanias). A manic episode is heralded by euphoric or irritable mood and at least four of the following symptoms: increased energy, activity, self-
esteem, or speed of thought; decreased sleep; poor judgment; and increased risk-taking. The energy, drive, and optimism of a hypomanic individual can seem captivating to others. About one-half of manic episodes are psychotic. The delusions of mania typically reflect grandiose or paranoid themes. Almost everyone who experiences manic episodes also will suffer from recurrent depressive episodes. Since there is a 50-50 chance that an episode of depression may precede the onset of the first mania, changes in diagnosis are common. The term bipolar affective disorder has largely replaced the former term manic-depression, although both names convey the cyclical nature of this illness. The classical presentation (which includes fullblown manic episodes) is known as type I disorder. The diagnosis of bipolar type II disorder is used when there are recurrent depressive episodes and at least one hypomania. A third diagnosis, cyclothymia, is used when neither hypomanias nor depressions have reached syndromal levels. For some, bipolar II or cyclothymic disorders subsequently evolve into the more severe type I form of illness. Several important variations of bipolar disorder are increasingly recognized. Mixed episode. A mixed episode is diagnosed when the symptoms of mania and depression coexist. Mixed episodes are diagnostically and therapeutically challenging, and their existence suggests that mania and depression are not caused by polaropposite abnormalities of brain functioning. Mixed episodes are often associated with sleep deprivation or drug and alcohol use. Rapid cycling. The term rapid cycling is used when there have been four or more episodes within a time frame of 1 year. Cases of ultrarapid cycling (mood cycles occur within hours or days) also have been reported. Rapid cycling is more common among women, and may be provoked by antidepressant therapies. Thyroid disease and the bipolar II subtype also increase the likelihood of rapid cycling. Seasonal affective disorder. A number of affective disorders follow a seasonal pattern. A pattern of recurrent fall/winter depressions (also known as seasonal affective disorder) has generated considerable interest because it may be treated with bright white (full-spectrum) light, which artificially lengthens the photoperiod. Disorders associated with general medical conditions. Both depression and mania can be caused by general medical illnesses and medications that affect brain function (such as antihypertensives, hormonal therapies, steroids, and stimulant drugs). The diagnosis “mood disorder associated with a general medical condition” is applied to these conditions. Recognition that the mood disorder is linked to another general medical illness can have important treatment implications. Pathophysiology Affective disorders have diverse biopsychosocial underpinnings that result, at least in part, in extreme or distorted responses of several neurobehavioral systems. The neurobehavioral systems of greatest
199
200
Affective disorders relevance regulate a person’s drives and pursuits, responses to acute stress, and capacity to dampen or quiet pain or distress. The system that regulates the daily rhythms of sleeping and waking and associated patterns of hormonal secretion likewise helps to ensure well-being. The coordinated function of these response and homeostatic systems conveys important and, from an evolutionary standpoint, ancient advantages for survival. Despite obvious overlap with how stress can adversely affect other primates and lower mammals, the affective disorders are distinctly human conditions. They may be thought of as one disadvantageous consequence of the complexity of affective, cognitive, and social functioning that distinguishes humans from their phylogenetic relatives. Although there is considerable evidence that affective disorders are heritable, vulnerability is unlikely to be caused by a single gene. Bipolar disorder is clearly more commonly inherited than major depressive disorders, and among the major depressive disorders the more severe, psychotic, and early-onset forms are more heritable than the remainder. Such risks are relative, however, and even identical twins have no more than a 50–60% risk of concordance of bipolar disorder. It is likely that some combination of genes conveys greater risk and, like an amplifier, distorts the neural signals evoked by stress and distress. See BEHAVIOR GENETICS; HUMAN GENETICS. Research permits several firm conclusions about brain neurochemistry in stress and depression. Acute stress affects the release of three vital brain monoamines—serotonin, norepinephrine, and dopamine—as well as glucocorticoids such as cortisol. Cortisol release is “driven” by the neuropeptide corticotrophin releasing hormone, which is released from neurons in both the hypothalamus and limbic cortex. Sustained unresolvable stress eventually depletes the neurotransmitters (cortisol levels remain high), inducing a behavioral state that has been called “learned helplessness.” In severe recurrent episodes of depression, especially those with psychotic features, cortisol elevations approach those seen in Cushing’s disease. Sustained elevations of cortisol have been shown to inhibit the growth of new neurons (neurogenesis) and ultimately may cause atrophy of some regions of the brain. See ADRENAL CORTEX HORMONE; BRAIN; STRESS (PSYCHOLOGY). As described earlier, both depression and mania are characterized by alterations in the sleep-wake cycle, including changes in nocturnal body temperature, growth hormone secretion, and the electroencephalographic patterns that define the stages of sleep. In depression, core body temperature tends to be elevated, growth hormone release is blunted, and there are increased awakenings. The brain’s patterns of electrical activity are altered, with decreased deep (slow-wave) sleep and increased rapid eye movement (REM) sleep in individuals with mania and depression. These changes suggest increased nocturnal arousal, perhaps because of decreased inhibitory serotoninergic tone. See SLEEP AND DREAMING. Positron emission tomography (PET) studies of brain function in depression document decreased
prefrontal cortex metabolism and increased activity of deeper, paralimbic structures such as the amygdala. The former disturbance may explain the cognitive abnormalities of severe depression, whereas the latter may reflect intense affective arousal. Mania, by virtue of the increased activity and characteristic changes in judgment and cooperativeness, has proved more difficult to study. In the prevailing view, it is a state of increased catecholamine activity, perhaps superimposed on a deficit of inhibitory serotoninergic input. Repeated studies of patients with bipolar affective disorder suggest that cerebral blood flow and glucose metabolism increase in prefrontal regions following a shift from depression to hypomania. Psychosocial and neurobiologic vulnerabilities, no doubt, intersect. For example, harsh early maltreatment, neglect, or other abuses can have lasting effects on both self-concept and brain responses to stress. Low self-esteem, in turn, can amplify the impact of a particular adverse event for an individual. Conversely, blunted hedonic capacity, perhaps mediated by low serotonin and catecholamine levels, decreases coping responses and the willingness to seek out new rewards. Similarly, chronic hyperarousal undoubtedly reinforces caution and more avoidant and dependent behaviors. Epidemiology The lifetime rates of affective disorders are increasing, with an earlier age of onset in more recent generations. The risks in the United States are for major depressive disorder, 10–15% of individuals; dysthymia, 3–5%; bipolar disorder type I, 1%; and bipolar disorder type II, 1–3%. Among these conditions, only bipolar disorder type I has a comparable incidence among men and women. For the remainder, female preponderance (1.5 or 2 to 1 ratios) is observed in most western cultures, which probably reflects both psychosocial and biological vulnerabilities. Minor depressive symptoms, including insomnia, constitute a major longitudinal risk factor. Other risk factors include divorced or single marital status and concomitant psychiatric, general medical, or substance-abuse disorders. The onset of major depression most often occurs in those in their late 20s to mid-30s; dysthymia and bipolar disorder typically begin about a decade earlier. However, no age group is immune to an affective disorder. Vulnerability is not strongly related to social class or race, although the affluent are more likely to receive treatment. Worldwide, the affective disorders are underrecognized and undertreated. In the United States, less than one-third of depressed people are receiving proper treatment, and only 60% of those with bipolar disorder are treated at any given time. Untreated depressive episodes tend to last months longer, increasing the risk of chronicity. Once chronicity is established, spontaneous remission is uncommon and episodes may last decades. About 70% of the first-onset major depressive episodes eventually become recurrent. After three or more episodes, the risk of further episodes (without
Affective disorders preventive treatment) exceeds 90%. This apparently increasing risk of recurrence is presumed to result from changes in brain stress responses referred to as “kindling.” The affective disorders, on average, cause more day-to-day disability than arthritis, diabetes, and most other common chronic illnesses. Depression’s effects on the quality of life essentially match that of congestive heart failure or chronic obstructive pulmonary disease. The World Health Organization ranks depression as the world’s fourth greatest health problem and estimates a second-place ranking in 2020. Bipolar disorder is ranked as the eighth greatest cause of illness burden. The lifetime risk of suicide is about 8–10% for the depressive disorders and 10–15% for bipolar disorder. Treatment Most episodes of dysthymia and major depressive disorder respond to treatment with either psychotherapy or antidepressant medication, either singly or in combination. Many experts now recommend the newer forms of psychotherapy, including cognitivebehavioral therapy and interpersonal therapy, because they have been better studied than more traditional psychoanalytic therapies and because they have been found to be as effective as medications. Nearly 30 antidepressant medications are available worldwide, with most falling into three classes: tricyclic antidepressants (TCAs), selective serotonin reuptake inhibitors (SSRIs), and monoamine oxidase reuptake inhibitors (MAOIs). Most classes of antidepressants enhance the efficiency of serotonin or norepinephrine neurotransmission. Antidepressants are not habit-forming and have no mood-elevating effects for nondepressed people. See MONOAMINE OXIDASE; NORADRENERGIC SYSTEM; SEROTONIN. The classes of antidepressants are not interchangeable, and a poor response or intolerable side effects with one medication do not preclude a good response to another. The SSRIs [in the United States, fluoxetine (Prozac), sertraline (Zoloft), paroxetine (Paxil), citalopram (Celexa), and escitalopram (Lexapro)], venlafaxine (Effexor), and bupropion (Wellbutrin) are most often prescribed as first-line treatments because of their tolerability, safety, and relative simplicity. It is hard to predict which depressed patient will respond best to which particular antidepressant. When treatment is effective, an improvement should be apparent within 4–6 weeks. An effective antidepressant should be maintained for 6–9 months to reduce the risk of relapse. Even longer prophylactic treatment is recommended for prevention of highly recurrent depressions. See PSYCHOPHARMACOLOGY; PSYCHOTHERAPY. There has been recent controversy about the possibility that antidepressants might paradoxically increase the likelihood of suicidality, especially early in the course of therapy. Although a cause-and-effect relationship has not been established, it is true that a small percentage of people who begin taking antidepressants (on the order of 1–3%) do feel worse
and report either the onset, or worsening, of suicidal thoughts or actions. Accordingly, regulatory authorities such as the U.S. Food and Drug Administration have required that the manufacturers of antidepressants include warnings about the possibility of treatment-emergent suicidality and the need for careful monitoring. A related issue concerns the use of antidepressants in children and teenagers. In brief, younger depressed people may be at greater risk for treatment-emergent suicidality and, in youth, the therapeutic benefits of antidepressant medications are not as well established. Acute manic episodes are usually treated with lithium salts, divalproex sodium, or one of the newer atypical antipsychotic medications, singly or in combination. The newer class of antipsychotic medication includes olanzapine (Zyprexa), risperidone (Risperdal), quetiapine (Seroquel), Ziprasidone (Geodon), and aripiperazole (Ability). The potent sedative benzodiazepines, such as lorazepam or clonazepam, also are used to control agitation and insomnia. Several other drugs, notably carbamazepine, are used when the primary agents are ineffective. Successful treatment of a manic episode should usually lead to preventive therapy. In addition to the medications used to treat mania, the anticonvulsant lamotrigine (Lamictal) is used for preventive care. Relapses that occur despite preventive treatment may require the use of more complex, multidrug regimens. Although psychotherapy does not have a major role in the acute treatment of mania, it may help people come to terms with their illness, cope more effectively with stress, or curb minor depressive episodes. Antidepressants also are used to treat “breakthrough” major depressive episodes. When pharmacotherapies are not effective, the oldest proven treatment of the affective disorders, electroconvulsive therapy (ECT), still provides a powerful alternative. Today, ECT is a highly modified and carefully monitored treatment that has little in common with its depictions in the movies. Nevertheless, confusion and transient amnesia are still common short-term side effects. Vigorous pharmacotherapy is needed after successful treatment to lessen the risk of relapse. See ELECTROCONVULSIVE THERAPY. Michael E. Thase Bibliography. American Psychiatric Association, Diagnostic and Statistical Manual of Mental Disorders, 4th ed. (DSM-IV), 1994; American Psychiatric Association, Practice guideline for the treatment of patients with major depressive disorder (revision), Amer. J. Psychiat., 157 (4 suppl.):1–45, 2000; F. K. Goodwin and K. R. Jamison, ManicDepressive Illness, Oxford University Press, 1990; R. M. A. Hirschfeld et al., Practice guideline for the treatment of patients with bipolar disorder (revision), Amer. J. Psychiat., 159(suppl. 4): 4:4–50, 2002; S. D. Hollon, M. E. Thase, and J. C. Markowitz, Treatment and prevention of depression, Psychol. Sci. Public Interest, 3(2):39–77, 2002; M. E. Thase, R. Jindal, and R. H. Howland, Biological aspects of depression, in I. H. Gotlib and C. L. Hammen (eds.), Handbook of Depression, pp. 192– 218, Guilford Press, New York, 2002.
201
202
Aflatoxin
Aflatoxin Any of a group of secondary metabolites produced by the common molds Aspergillus flavus and A. parasiticus that cause a toxic response in vertebrates when introduced in low concentration by a natural route. The group constitutes a type of mycotoxin. Discovered in 1960 after a massive poisoning of turkey poults fed moldy peanut meal, aflatoxin has become the focus of a highly interdisciplinary research field involving chemists, mycologists, veterinarians, agriculturalists, toxicologists, and other basic and applied scientists. See TOXIN. Chemistry. The naturally occurring aflatoxins are identified in physicochemical assays as intensely blue (aflatoxins B1 and B2) or blue-green (aflatoxins G1 and G2) fluorescent compounds under long-wave ultraviolet light. The common structural feature of the four major aflatoxins is a dihydrodifurano or tetrahydrodifurano group fused to a substituted coumarin group (see illus.). The relative proportions of the four major aflatoxins synthesized by Aspergillus reflect the genetic constitution of the producing strain and the parameters associated with fungal growth. In addition, derivative aflatoxins are produced as metabolic or environmental products. Aflatoxins are formed through a polyketide pathway involving a series of enzymatically catalyzed reactions. In laboratory cultures, aflatoxins are biosynthesized after active growth has ceased, as is typical for secondary metabolites. By using blocked mutants and metabolic inhibitors, many of the intermediates have been identified as brightly colored anthraquinones. Working with these colored metabolites has made studies of aflatoxin biosynthesis into a model system for polyketide research and has also facilitated molecular genetics. Several genes encoding enzymes that catalyze steps in the pathway or affect more global regulatory mechanisms have been cloned by using complementation studies in which transformants are selected based on pigmentation. Biological effects. Aflatoxins are potent molecules with many biological effects. They are toxigenic, carcinogenic, mutagenic, and teratogenic in various animal species. Aflatoxin B1 is usually the most abundant naturally occurring member of the family, and most studies on the pharmacological activity of aflatoxin have been conducted with this congener. Interspecific and interstrain variation is the rule for both acute and chronic effects. Trout and ducklings are among the most sensitive species known. The toxic properties manifest themselves in different ways depending on the test system used, the dose given, and the duration of exposure. In addition, age, nutritional status, and sex affect the toxicological response; in general, the young of a species are the most sensitive. In veterinary medicine, aflatoxicosis is often difficult to detect. Inhibited growth, weight loss, and immune response breakdown are observed, and subtle synergies with both infectious and nutritional diseases complicate the diagnostic profile. The main target organ for aflatoxin is the liver. Aflatoxin B1 is
O
O
O
O
OCH3
O
(a) O
O
O
O
OCH3
O
(b) O
O
O
O
O
OCH3
O
(c) O
O
O
O
O
O
OCH3
(d) Structures of major naturally occurring aflatoxins. (a) B1. (b) B2. (c) G1. (d) G2.
the most potent hepatocarcinogenic agent known, although the liver by no means is the only organ susceptible to aflatoxin carcinogenesis. Epidemiological studies in Africa and Asia link ingestion of mold-contaminated food with incidence of primary liver cancer. Moreover, independent studies on liver cancer from southern Africa and Qidong, China, revealed guanine to thymine transversions in a tumor suppressor gene. Identification of this molecular hot spot constitutes some of the most convincing evidence that a natural product in a nonindustrial setting can induce a specific form of cancer. Aflatoxin is listed as a probable human carcinogen by the International Agency for Research on Cancer. See BOTULISM; LIVER DISORDERS; MUTAGENS AND CARCINOGENS; PLANT PATHOLOGY. Economics. Aflatoxins are a major agricultural problem. Contamination can occur in the field, during harvest, or in storage and processing. Corn, rice, cottonseed, and peanuts are the major crops
Africa regularly displaying high levels of aflatoxin contamination. Since A. flavus and A. parasiticus are nearly ubiquitous in the natural environment, numerous other grain, legume, nut, and spice crops, as well as coffee and cocoa, have been reported to contain aflatoxins. Given the potential of aflatoxins as human carcinogens and their known activity as toxins in animal feeds, many international regulatory agencies monitor aflatoxin levels in susceptible crops. Contaminated agricultural commodities typically contain low amounts of aflatoxins (in the parts per billion or parts per million range) so numerous methods have been devised for extraction, concentration, and purification. Since agricultural commodities are diverse, no single extraction procedure can be applied uniformly to every product. Most of the original analytical methods employ thin-layer chromatography or high-performance liquid chromatography. More recently, immunological assays have become popular because they require fewer extractions and less complicated equipment. Prevention is the main line of defense against aflatoxins entering the food chain. Moisture, temperature, and composition of the substrate are the chief factors affecting fungal growth and toxin production. In the field, insect damage is often involved. Good agronomic techniques that minimize damage to crops, storage under conditions unfavorable to mold growth, and screening of crops in order to exclude those which have been contaminated are the major methods for protecting the food supply. In some cases, portions of crops may be salvaged by mechanical separation. Detoxification is a last line of defense. Several commercially feasible methods of ammoniation have been developed for reducing levels of aflatoxin contamination in animal feeds. Safety. Laboratory workers handling aflatoxins, especially in purified forms, should maintain strict safety precautions. Aflatoxin-contaminated foods should not be eaten or fed to domesticated animals. See AGRONOMY; MYCOTOXIN. J. W. Bennett Bibliography. J. W. Bennett and M. A. Klich (eds.), Aspergillus: Biology and Industrial Applications, 1992; R. J. Cole and R. H. Cox, Handbook of Toxic Fungal Metabolites, 1981; J. E. Smith and M. S. Moss, Mycotoxins: Formation Analysis, and Significance, 1985.
Africa A continent that straddles the Equator, extending between 37◦N and 35◦S. It is the second largest continent, exceeded by Eurasia. The area, shared by 55 countries, is 11,700,000 mi2 (30,300,00 km2), approximately 20% of the world’s total land area. Despite its large area, it has a simple geological structure, a compact shape with a smooth outline, and a symmetrical distribution of climate and vegetation. Structure and geology. The structural makeup is partially explained by the continental drift theory, in which all of the Earth’s continents were previously assembled into one supercontinent called Pangaea
(or Pangea). Over time, Pangaea broke into two landmasses, Laurasia and Gondwana, and subsequently into several continents. Africa occupied the core of Pangaea, and as neighboring landmasses drifted away, the present structure and configuration of the continent evolved. See CONTINENTAL DRIFT; CONTINENTS, EVOLUTION OF; PLATE TECTONICS; SUPERCONTINENT. Three major rock formations can be identified. The first is a massive crystalline shield (Precambrian) that underlies most of the continent, with occasional outcrops in about one-third of the area. In some places, these rocks have been subject to metamorphic processes, resulting in a rich variety of mineral resources. The second set of significant rocks is the Karoo series in southern and eastern Africa. These are sedimentary rocks varying in age from the Carboniferous to the lower Jurassic Period. They are presently vital sources of coal. A third set of rock formations comprises marine sediments stretching in a discontinuous belt from the northwest to the Gulf of Guinea. These were formed during the Cretaceous Period, when parts of the continent were submerged in water. The sediments deposited at the time contain hydrocarbons that now account for oil production in Libya, Algeria, Nigeria, Gabon, and Angola. See CARBONIFEROUS; CRETACEOUS; JURASSIC; PRECAMBRIAN. Africa has few inlets or natural harbors and a small number of offshore islands that are largely volcanic in origin. Madagascar is the largest island, with an area of 250,000 mi2 (650,000 km2). It is separated from the mainland by the Mozambique Channel, which is about 250 mi (400 km) wide. Close to Madagascar is the Mascarene, a volcanic ridge submerged in the Indian Ocean. Its highest peaks are exposed above the sea surface to form the islands of Seychelles, Mauritius, and R´eunion. North of Madagascar, there is another group of volcanic islands, the Comoros. The largest of these is Njazide (Grand Comore), with an active volcano, Mount Kartala (7677 ft; 2340 m). See INDIAN OCEAN. Surface features. Africa is primarily a high interior plateau bounded by steep escarpments. These features show evidence of the giant faults created during the drift of neighboring continents. The surface of the plateau ranges from 4000–5000 ft (1200– 1500 m) in the south to about 1000 ft (300 m) in the Sahara. These differences in elevation are particularly apparent in the Great Escarpment region in southern Africa, where the land suddenly drops from 5000 ft (1500 m) to a narrow coastal belt. Although most of the continent is classified as plateau, not all of its surface is flat. Rather, most of its physiographic features have been differentially shaped by processes such as folding, faulting, volcanism, erosion, and deposition. See ESCAPEMENT; FAULT AND FAULT STRUCTURES; PLATEAU. Basin surfaces. In several areas the plateau has sagged under the weight of accumulated riverine sediments to form large structural basins. These include the Sudan, Chad, and Djouf basins in the north, Congo (Zaire) basin in central Africa, and the Kalahari basin
203
204
Africa in the south. These basins play a significant role in the hydrography of the continent. See BASIN; HYDROGRAPHY. Rift valley system. The old crustal surfaces in the eastern part of the continent have been disrupted by a major structural feature known as the East African rift valley system that consists of a complex set of troughs formed by the downward displacement of land blocks against one another. The origin of the rift valley system can be traced to movements associated with separation of the Arabian plate from the African plate. The resulting landscape emanates from the Red Sea and passes through the Ethiopian highlands, along the course of the Awash river and a number of lakes (including Zwai, Shala, Abaya, and Turkana) into western Kenya. From there, it cuts across the Kenyan highlands and a series of lakes from Lake Banyo to Lake Magadi. It then enters Tanzania via Lake Natron, after which it loses its troughlike appearance as a result of the accumulation of sediments and volcanic ash. The system reappears in its original form as it passes through Lake Malawi and the Shire river into the lower Zambezi river and finally enters the Indian Ocean. See RED SEA. The western arm of the rift valley starts north of Lake Albert (Lake Mobutu) in Uganda and passes through Lakes Edward (Amin), Kivu, and Tanganyika into Lake Nyasa, with a trench running southwest as the Luangwa valley of Zambia. Lake Tanganyika is the continent’s deepest lake (5100 ft; 1500 m) with a bottom that is roughly 0.5 mi (800 m) below sea level. The rift valley system is one of the most striking features of the African landscape. Sliding blocks have created wide valleys 20–50 mi (30–80 km) wide bounded by steep walls of variable depth and height. Within the eastern and western branches of the system, there is a large but shallow depression occupied by Lake Victoria. Between Lake Edward and Lake Albert on the western extension is the icecapped Ruwenzori range, a block mountain, that rises over 16,732 ft (5100 m). These areas, along with other parts of the continent, occasionally experience earthquakes; the largest earthquake in the twentieth century occurred near the southern end of Lake Tanganyika. See RIFT VALLEY. Volcanic landforms. Several volcanic features are associated with the rift valley system. The most extensive of these are the great basalt highlands that bound either side of the rift system in Ethiopia. These mountains rise over 10,000 ft (3000 m), with the highest peak, Ras Dashan, reaching 15,158 ft (4500 m). There are also several volcanic cones, including the most renowned at Mount Elgon (14,175 ft; 4321 m); Mount Kenya (17,040 ft; 5194 m); and Mount Kilimanjaro, reaching its highest point at Mount Kibo (19,320 ft; 5889 m). Mounts Kenya and Kilimanjaro are permanently snowcapped. See BASALT. There are spectacular volcanic landforms in other parts of the continent such as the Ahaggar (Hoggar) mountains in central Sahara, which has over 300 lava plug domes rising above 1000 ft (305 m). There are also huge craters or calderas scattered in various
places such as Trou au Natron in the Tibesti mountains and the Ngorongoro in Tanzania. Present-day active volcanoes include Oldoinyo Lengai in Tanzania; Nyamlagira and Nyriragongo in the Virunga range on the Congo (Zaire)–Uganda border; Erta Ale in the Afar Triangle of Ethiopia and on several offshore islands such as San Juan on Las Palmas islands; Fogo on Cape Verde islands; Kartala on Comoros islands; and Piton de la Fournaise on R´eunion island. The largest active volcano is Mount Cameroon (13,350 ft; 4069 m). Although it lies on the mainland, it is part of a chain of volcanoes extending into the Atlantic Ocean. The islands of Fernando Po, Principe, and S˜ao Tom´e are also part of this system. See OCEANIC ISLANDS; VOLCANO. Fold mountains. The northwestern and southern margins of the continent have been subject to folding and subsequent erosion, resulting in a complex chain of mountains and valleys. The Atlas ranges in the northwest extend in a belt up to 200 mi (320 km) wide from southern Morocco to Tunisia, a distance of 1400 mi (2250 km). The system consists of a series of parallel mountains built by folding and uplifting processes during the Tertiary Period. Geologically they are an extension of the Alpine system of Europe, now separated by the Mediterranean Sea. The two principal chains are found along the coast and in the interior. The coastal chain starts from Tangiers in Morocco as the Rif, extends eastward into Algeria as the Tell, and ends in Tunisia. The interior chain starts west of Morocco as the High Atlas and Anti-Atlas. It continues eastward into Algeria as the Sahara Atlas and terminates in Tunisia. Between the two chains lies the Plateau of Shott, a salty shallow basin with internal drainage. Another plateau, the Moroccan Meseta, is found in the western extreme, sloping toward the Atlantic Ocean. See TERTIARY. The Cape ranges in southern Africa are less complex, consisting of folded sedimentary rocks that correspond in age and topography to Appalachia in the eastern United States. The mountains run in a northsouth direction along the west coast for about 150 mi (240 km) and then make right angles to run east-west for another 600 mi (900 km). See OROGENY; SEDIMENTARY ROCKS. Climatic factors. Since the Equator transects the continent, the climatic conditions in the Northern Hemisphere are mirrored in the Southern Hemisphere. Nearly three-quarters of the continent lies within the tropics and therefore has high temperatures throughout the year. Frost is uncommon except in mountainous areas or some desert areas where nighttime temperatures occasionally drop below freezing. These desert areas also record some of the world’s highest daytime temperatures, including an unconfirmed record of 136.4◦F (58◦C) at Azizia, Tripoli. See EQUATOR. Africa is largely under the influence of subtropical high-pressure belts separated by the equatorial lowpressure zone known as the Inter-Tropical Convergence Zone (ITCZ). The position of this zone across the tropics shifts with the seasonal movement of the Sun. In July, it moves to North Africa along the
Africa margins of the Sahara. In January, it moves south, taking up a clearly defined position in West Africa. Its position in eastern and southern Africa is not well defined because of the elevational variations there, but it tends to slope along the eastern margins of the Democratic Republic of Congo (formerly Zaire) and eastward through Mozambique. The Inter-Tropical Convergence Zone plays a fundamental role in the circulation of the airstreams and subsequent climatic conditions within the continent. The associated air masses are the tropical maritime and the tropical continental air with notable differences in humidity due to their source areas. The tropical maritime air is warm and moist, originating from the Atlantic Ocean. The tropical continental air is warm and dry because it develops over the desert areas. In January, following the southern movement of the Inter-Tropical Convergence Zone, most of the Northern Hemisphere falls under the influence of tropical continental air, bringing in stable dry weather and dusty haze known as the harmattan in West Africa. In July, the shifting Inter-Tropical Convergence Zone and the tropical maritime air take over the region, bringing in heavy rainfall. See AIR MASS; CLIMATOLOGY; TROPICAL METEOROLOGY. The wind reversals associated with the InterTropical Convergence Zone produce distinct monsoonal seasons along the coastal margins of eastern and western Africa. The conditions observed along the eastern margins (including the island of Madagascar) are part of the larger Asian monsoonal system. During the wet summer monsoon, the northern shift of the Inter-Tropical Convergence Zone results in moist onshore winds from the warm tropical Indian Ocean. The atmospheric conditions are reversed in January, when dry offshore winds prevail following the southern shift of the Inter-Tropical Convergence Zone. See MONSOON METEOROLOGY. The West African monsoon is smaller in magnitude than the Asian monsoonal system observed in eastern Africa. It dominates within 400 mi (650 km) of the coastline, particularly along the Gulf of Guinea. The causal factors are similar, however: the shifting wind belts produce alternate wet and humid conditions in the summer and dry and dusty conditions in winter. Finally, the convergence of surface winds along the Inter-Tropical Convergence Zone is also responsible for the Easterly Wave, which occurs when the Inter-Tropical Convergence Zone is displaced away from the Equator, usually between 5◦ and 20◦ latitude. The weak wave is an area of low pressure that is propelled westward within the trade wind belt. Given the appropriate surface and weather conditions, these systems can develop into powerful, wellorganized tropical disturbances that eventually grow into hurricanes. Most of the tropical storms and hurricanes in the northern Atlantic are linked to the Easterly Wave originating off the western coast of Africa. However, not all storms grow into hurricanes. Hurricanes typically develop where surface water temperatures are warm, 79◦F (26◦C) or higher, over a
large area with sustained winds exceeding 74 mi/h (119 km/h). Regional distribution of climate, vegetation, and soils. Africa can be classified into broad regions based on the climatic conditions and their associated vegetation and soil types. The tropical rainforest climate starts at the Equator and extends toward western Africa. The region has rainfall up to 200 in. (500 cm) per year and continuously high temperatures averaging 79◦F (26◦C). The eastern equatorial region does not experience these conditions because of the highlands and the presence of strong seasonal winds that originate from southern Asia. The tropical rainforest vegetation consists of tall broadleaf trees that are categorized into different layers. The topmost layer, or emergents, consists of tall, straight trees that average 165 ft (50 m) in height. The middle layer, with trees of 80–155 ft (25–35 m), forms a dense interlacing canopy that allows little sunlight to penetrate to the forest floor. The lowest layer consists of young trees that average 50 ft (15 m) in height. The forest floor is typically damp and relatively open except for dead leaves, fallen branches, logs, and woody vines that connect the the tree trunks to the canopy. The areal extent of the African rainforest region (originally 18%) has dwindled to less than 7% as a result of high rates of deforestation. Despite these reductions, the region is still one of the most diverse ecological zones in the continent. It has not only supported the timber industry in several countries but also served as an excellent reservoir for wildlife and medicines for traditional herbalists. Estimates from the International Union for Conservation of Nature show that the forested regions along the West African coast and within the central African basin contain more than 8000 plant species, including a large variety of hard and soft woods. This region also has about 84% of Africa’s primate species and 68% of its passerine birds. Most of the mammals are arboreal, but there are several ground-dwelling animals such as rodents, wild pigs, and deer. Insects generally outnumber all other animal species combined. Typical are lateritic soils or oxisols with limited plant nutrients due to accelerated rates of decomposition and severe leaching. See RAINFOREST. The tropical rainforest gradually merges northward and southward into the tropical savanna. This savanna encompasses a more extensive area and is best characterized by its distinct wet and dry seasons. Annual rainfall is lower (40 in.; 100 cm), and the mean annual temperature is higher than the equatorial zone. Vegetation includes savanna woodlands and shrubs interspersed with coarse tall grasses. Unlike the forested zones, there is no continuous tree canopy. The woodlands consist of open stands of trees that are at least 16 ft (5 m) tall. The two major zones are the Guinea savanna region in West Africa and the Zambezian region in eastern Angola, the Democratic Republic of Congo, and Zambia. Dominant tree species include the African fan palm (Borassus aethiopicum), silk cotton tree (Bombax petandrum), boabab (Adansonia digitata), African
205
206
Africa rubber (Landolphia spp.), and mopane and miombo tree species. In moving farther north or south of the Equator, tree size and overall density gradually decrease, with an increasing coverage of herbaceous plants and grassland. Extensive savanna grasslands are found along the Sudanian zone of West Africa, within the Zambezian region and the Somalia-Masai plains. Large areas such as the Serengeti plains in the Somalia-Masai plains are home to a diverse range of wild animals. The soils are well developed and rich in nutrients because they are not subject to severe erosion. However, in some areas, especially where the topsoil has been removed, they have developed into thick crusts or pans during the dry seasons. The volcanic regions in eastern Africa are an exception, having productive soils. The tropical steppe forms a transition zone between the humid areas and the deserts. This includes the area bordering the south of the Sahara that is known as the Sahel, the margins of the Kalahari basin, and the Karoo grasslands in the south. Average temperatures range from a daily high of 87◦F (31◦C) in August to 107◦F (42◦C) in April. Annual rainfall is highly variable, averaging 10–20 in. (25–50 cm). Inadequate rainfall coupled with poor soil chemistry (such as lack of nutrients) limits plant growth and species diversity. Nonetheless, this region has been used extensively for pastoral farming and for growing staples such as millet and sorghum. It is the most vulnerable ecological zone in the continent. In particular, the Sahelian zone, stretching from Mauritania and Senegal in the west to southern Sudan and Ethiopia in the east, has gained worldwide attention as a result of a series of devastating droughts that have affected millions of lives since the 1970s. The underlying causes for these problems are complex and intricately tied to both physical environmental and human-made factors. Climatically, the most significant cause of desertification is the variability and marked reduction in rainfall beginning in the 1970s. This has been partly attributed to global warming trends. Specifically, higher sea surface temperatures over the Atlantic tend to hinder the northward penetration of the southwest monsoonal winds, thus reducing the likelihood for precipitation. An alternative explanation links the observed drought conditions to dust storms that frequently occur in the Sahara. These tend to heat up the atmosphere and prevent the upward flow of air needed for precipitation to occur. See DESERTIFICATION. Aside from physical factors, the desertlike conditions in the Sahel are partly the result of specific human activities that degrade land. These activities include the removal of stabilizing natural vegetation for fuel wood and other land uses, overcultivation of the marginal areas, and overgrazing. Also contributing are a host of socioeconomic and political problems in the affected countries. As one moves farther away from these marginal areas, the steppe trends into the desert landscape in the southwestern and northern parts of the conti-
nent. The Namib desert borders the Atlantic Ocean in southern Africa. Inland, it merges into the Kalahari, more accurately described as a semidesert instead of a true desert. In the north the world’s largest desert, the Sahara, covers about 30% of the continent. The desert consists of mostly flat, dry plains and plateau with a few mountains such as the Ahaggar, the Air, and the Tibesti. Unlike most deserts around the world, more than 20% of the desert surface is covered with ergs (rolling plains of sand) such as the Libyan Erg (close to 200,000 mi2 or 80,940 km2) and the eastern and western ergs of Algeria (each about 20,000 mi2 or 8094 km2). Given this vast areal extent, the landscape takes on an undulating appearance with ripples of sand occasionally interrupted by small groups of dunes. In other areas, the desert landscape is characterized by regs (plains covered by pebbles and boulders), hammadas (polished rock surfaces), and basins of interior drainage (some containing salt pans). See DESERT. Average annual rainfall in the desert is consistently below 10 in. (25 cm) with extreme variability from year to year. The relative humidity is very low (under 5% and rarely above 20%). There are wide annual temperature ranges and extreme diurnal ranges. The most common weather occurrences are sandstorms, particularly in the winter and early spring. Depending on the velocity and the volume of sand being transported, these storms can be disastrous for crops, livestock, and the inhabitants of this region. Extensive studies conducted in the last few decades have focused on the characteristics of these storms in relation to the materials deposited elsewhere (such as the Americas, Asia, and Europe). Specifically, airborne dust from the Sahara negatively impacts the ecology and climate within the region itself but provides valuable loess deposits in areas westward across the Atlantic Ocean, northward over the Mediterranean Sea, and eastward across the Red Sea. See LOESS. Land-use patterns are extremely variable in the deserts. They range from extensive nomadic grazing of livestock to intensive cultivation around oases. These activities, however, depend on the local variation in climate, soil moisture, and the availability of water through traditional irrigation systems. The most important staple is the date palm. Fruits, vegetables, and cereals are also grown. The northwestern and southwestern margins of the continent have a mediterranean climate, with dry summers and cool wet winters. The latter are caused by the polar air mass which brings relatively low temperatures and rainfall. Annual rainfall is about 30 in. (75 cm) or less. Mostly drought-resistant trees and scrubs are found here. The soils are variable, depending on the underlying bedrock, but can be generally classified as acidic and less fertile. See PLANT GEOGRAPHY; SOIL. Drainage systems. The structural evolution of the continent has much to do with the drainage patterns. Originally, most of the rivers did not drain into the oceans, and many flowed into the large structural basins of the continent. However, as the continental
African horsesickness drift occurred and coasts became more defined, the rivers were forced to change courses, and flow over the escarpments in order to reach the sea. Several outlets were formed, including deep canyons, waterfalls, cataracts, and rapids as the rivers carved out new drainage patterns across the landscape. Most of the rivers continue to flow through or receive some of their drainage from the basins, but about 48% of them now have a direct access into the surrounding oceans. The major rivers are the Nile, Congo (Zaire), Niger, and Zambezi. See RIVER. Nile. This is the only perennial river that enters the Mediterranean Sea from Africa. With an overall length of 4157 mi (6690 km), it is the longest river. The water originates from three major sources. About 56% comes from the Blue Nile, 22% from the Atbara River, and the rest from the White Nile. At Khartoum the White Nile joins the Blue Nile, and the plain between the two rivers, the Gezira, is the site of a large irrigation project. The Nile is controlled by several dams. The largest is the Aswan High Dam, which created Lake Nasser, over 300 mi (480 km) long. See DAM; MEDITERRANEAN SEA. Congo. The Congo (Zaire), which is 3000 mi (4800 km) long, is Africa’s largest river in terms of discharge. Its flow is stable throughout the year because it receives water from tributaries that lie on both sides of the Equator. The river originates near Lake Tanganyika and descends over a series of rapids and falls on its way down the Great Escarpment to the Atlantic Ocean. The Inga Dam, built between Kinshasa and Matadi, was completed in 1972. Large sections of the river are also navigable, particularly its middle course, which provides a 1085-mi (1746-km) stretch from Kinshasa to Kisangani. See ATLANTIC OCEAN. Niger. The Niger rises in the Guinean highlands, an area of heavy rainfall in West Africa. After a total course of 2600 mi (4160 km), it enters the Atlantic through a large delta in Nigeria. It was dammed in 1969 at Kainji, Nigeria, creating a large artificial lake. Other large rivers in this region include the Senegal and the Volta. See DELTA. Zambezi. The Zambezi begins in Angola, and its 2200-mi (3540-km) course drains most of the area south of the Congo system into the Indian Ocean. There are several obstructions in its middle course, including Victoria Falls. Here the Zambezi plunges 350 ft (107 m) from the inland Kalahari basin onto the edge of hard resistant rocks, cutting deep canyons in its path. A dam was built downstream from the falls in 1959, creating Lake Kariba. Also in Mozambique, the Caborra Bassa dam was completed in 1976 at the point where the river leaves the Kebrabasa gorge on its way to coastal lowlands. Other major rivers within southern Africa include the Limpopo and the Orange. Fauna. The tremendous diversity in wildlife continues to be one of the primary attractions of this continent. Africa is one of the few remaining places where one can view game fauna in a natural setting. There is a tremendous diversity in species, including birds, reptiles, and large mammals. Wildlife are
concentrated in central and eastern Africa because of the different types of vegetation which provide a wide range of habitats. Until recently, these regions were the least disrupted areas. Human population densities were relatively low except within highland and fertile volcanic areas. However, the situation is changing rapidly as the animals face constant threats, including disease, frequent drought, poaching, and population encroachment. Efforts to save them include the development of game reserves, bans on hunting or poaching, controlling exports of wildlife (or their products), and wildlife management. See CONSERVATION OF RESOURCES. Population and environmental hazards. Africa is not a densely populated continent. With an estimated population of 743,000,000, its average density is 64 per square mile (26 per square kilometer). However, some areas have large concentrations, including the Nile valley, the coastal areas of northern and western Africa, the highland and volcanic regions of eastern Africa, and parts of southern Africa. These are mostly areas of economic or political significance. In several regions, the human population is exposed to a variety of hazards such as riverine flooding, drought, disease, and pests. These hazards are largely meteorological in origin, such as the long periods of drought in the Sahel and eastern and southern Africa. The climate also provides a habitat for harmful insects and other vectors of endemic diseases such as malaria, yellow fever, trypanosomiasis, onchocerciasis, and bilharzia. Crops are occasionally attacked by swarms of locusts and other pests. See CONTINENT; EPIDEMIOLOGY. Florence Lansana Margai Bibliography. S. Aryeetey-Attoh (ed.), Geography of Sub-Saharan Africa, Prentice Hall, 1997; A. Grainger, The Threatening Desert: Controlling Desertification, Guersney Press, 1982; I. L. Griffiths, An Atlas of African Affairs, University of Witwatersrand Press, 1994; J. Gritzner, The West African Sahel, University of Chicago, Committee on Geographic Series, 1988; International Union for Conservation of Nature, The Conservation Atlas of Tropical Forests in Africa, Simon and Schuster, 1992; L. A. Lewis and L. Berry, African Environments and Resources, 1988; A. B. Mountjoy and D. Hilling, Africa: Geography and Development, Barnes and Noble, 1987; S. E. Nicholson, Climatic variations in the Sahel and other African regions during the past five centuries, J. Arid Environ., 1:3–24, 1987; J. M. Pritchard, Landform and Landscape in Africa, Arnold, 1979; P. Richards, The environmental factor in African studies, Prog. Human Geog., 4:589–600, 1980; C. Stager, Africa’s Great Rift, Nat. Geog., pp. 2–14, May 1990.
African horsesickness A highly fatal insect-borne viral disease of horses and mules, and a mild subclinical disease in donkeys and zebras. It normally occurs in sub-Saharan Africa but occasionally spreads to North Africa, the Iberian Peninsula, and Asia Minor.
207
208
Afterburner Etiologic agent. The African horsesickness virus is an orbivirus (family Reoviridae) measuring 68– 70 nanometers in diameter. The outer layer of the double-layered protein shell is ill defined and diffuse and is formed by two polypeptides. The highly structured core consists of five structural proteins arranged in icosahedral symmetry. The viral genome is composed of 10 double-stranded ribonucleic acid (RNA) segments (genes) ranging in size from 240,000 to 2,530,000 daltons. Nine distinct serotypes which can be distinguished by neutralization tests are known. The virus can be cultivated in various cell cultures, in the brains of newborn mice, and in embryonated hen eggs by intravascular inoculation. See ANIMAL VIRUS. Transmission. African horsesickness is a noncontagious disease that can readily be transmitted by the injection of infective blood or organ suspensions. In nature, the virus is biologically transmitted by midges of the genus Culicoides, such as C. imicola. The disease has a seasonal incidence in temperate regions (late summer to autumn), and its prevalence is influenced by climatic conditions favoring insect breeding (for example, warm, moist conditions in low-lying areas). Mechanical transmission by large biting flies is possible, but plays a much smaller role than biological transmission in the epidemiology of this disease. Clinical signs. These indications result from selective replication of the virus in vascular endothelial cells, leading to increased vascular permeability in specific organs or regions of the body. Clinically, four forms of the disease are recognized. (1) In the peracute, pulmonary form, an incubation period of 3–5 days is followed by a febrile reaction and signs of respiratory distress. The breathing accelerates (75 breaths per minute or more), and the horse stands with its neck extended and the nostrils fully dilated. Profuse sweating is common, and terminally spasmodic coughing may be observed with frothy fluid exuding from the nostrils. Death usually occurs within hours after the first respiratory signs are observed. (2) In the subacute, edematous or cardiac form, the incubation period is about 7– 14 days. The febrile reaction is followed by characteristic swellings of the hollows above the eyes, the eyelids, the lower lip, and the neck region. Swelling of the limbs is never observed. Some animals may show signs of colic, and death occurs 4–8 days after the onset of the febrile reaction. (3) Most horses show the acute or mixed form, which represents a combination of the pulmonary and edematous forms. (4) African donkeys, zebras, and partially immune horses may exhibit the fever form, where a mild febrile reaction is the only clinical sign. The mortality rate in fully susceptible horses varies from about 70 to 95% and is about 50% in mules. In European and Asian donkeys, the fatality rate may be about 10% while African donkeys and zebras never succumb to the disease. Pathology. The gross autopsy findings depend largely on the clinical form of disease manifested by the animal before death. In the pulmonary form,
large amounts of fluid are found in the lung tissue or in the thoracic cavity, and death results from suffocation. Other lesions commonly observed are intense inflammation of the glandular fundus of the stomach (which can easily be confused with poisoning) and hemorrhage or swelling of the large and small intestines. In the subacute form, there is extensive infiltration of subcutaneous and intermuscular tissues with yellow gelatinous fluid. In mild cases just the head and neck are involved, but in more severe cases the edema also involves the brisket and shoulders. Accumulation of fluid in the pericardial sac is a constant and very significant finding. Extensive hemorrhages are found on the surface of the heart as well as in the left ventricle. The lymph nodes are generally enlarged and edematous, and lesions in the gastrointestinal tract are largely similar to those found in the pulmonary form except that submucosal edema of the large intestine tends to be more conspicuous. In the mixed form the lesions represent a combination of those observed in the pulmonary and edematous forms. Treatment and control. There is no specific treatment for this disease. Infected animals should be given complete rest as the slightest exertion may result in death. Stabling of horses at night during the African horsesickness season reduces exposure to insect vectors and hence the risk of disease transmission. Prophylactic vaccination is the most practical and effective control measure. In epidemic situations outside Africa, the causal virus should be serotyped as soon as possible, allowing the use of a monovalent vaccine. However, in endemic regions it is imperative to use a polyvalent vaccine, which should render protection against all nine serotypes of African horsesickness virus. See DISEASE; EPIDEMIC; VACCINATION. Baltus Erasmus
Afterburner A device in a turbojet aircraft engine, between turbine and nozzle, which improves thrust by the burning of additional fuel (Fig. 1). To develop additional
combustion chamber inlet diffuser
flame holder tailcone
< < fuel nozzles <
1
Ma>1
c d, 0
0
Macr
Madd
1.0
Ma∞ Fig. 3. Variation of the profile drag coefficient with the free-stream Mach number Ma∞, illustrating the critical Mach number Macr and drag-divergence Mach numbers Madd and showing the large drag rise near Mach 1. Inset shows the flow field at velocities above that of the drag-divergence Mach number, with large regions of locally supersonic flow (Ma > 1).
Airframe for a given airfoil at a given angle of attack, that freestream Mach number at which sonic flow is first encountered at a point on the airfoil is the critical Mach number. When the free-stream Mach number is increased slightly above the critical Mach number, whole regions of locally supersonic flow occur over the airfoil (Fig. 3). These supersonic regions are terminated by a shock wave. If the shock wave is strong enough, its interaction with the boundary layer separates the flow at the shock impingement point. Both the shock wave itself and the flow separation caused by the shock contribute to a massive increase in drag and a decrease in lift. The onset of these phenomena is best described by considering the measurement of the drag coefficient of a given airfoil at a fixed angle of attack in a wind tunnel, as a function of Mach number (Fig. 3). The drag coefficient remains relatively constant as the Mach number is increased from an initial low value all the way up to the critical Mach number. As the Mach number is increased slightly above the critical value, a finite region of supersonic flow appears on the airfoil. With a further small increase in Mach number, a point is encountered where the drag coefficient suddenly starts to increase. The value of the Mach number at which this sudden increase in drag starts is defined as the drag-divergence Mach number. Beyond the drag-divergence Mach number, the drag coefficient can become very large, typically increasing by a factor of 10 or more. The peak value of the drag coefficient was so daunting to aeronautical engineers in the 1930s and early 1940s that it gave rise to the concept of the sound barrier, that is, that airplanes could not fly faster than sound (faster than Mach 1). Of course, there is no such barrier to supersonic flight. The supersonic transport Concorde cruises at Mach 2.2, and some high-speed military airplanes can fly at Mach 3 and above. The advent of jet propulsion in the 1940s ushered in the era of high-speed flight and focused attention on designing airfoils with high critical Mach numbers, thus trying to delay the transonic drag divergence to higher subsonic Mach numbers. One well-known approach was to use thin airfoils on high-speed aircraft; everything else being equal, a thin airfoil has a higher critical Mach number than a thicker airfoil. In the 1960s, a new airfoil design philosophy was developed, the supercritical airfoil. For the same free-stream Mach number, the supercritical airfoil is shaped so as to produce a smaller region of supersonic flow with a lower local supersonic Mach number than the more conventional airfoil shape. Because of this behavior, the supercritical airfoil has a larger gap between the critical and the drag-divergence Mach numbers; that is, it can penetrate farther into that part of the free-stream Mach number regime above the critical Mach number before encountering the drag-divergence phenomena. See SUPERCRITICAL WING; TRANSONIC FLIGHT. Supersonic flow. When the free-stream Mach number is greater than 1, the physical aspects of the flow over an airfoil are quite different from those for a freestream Mach number less than 1. In supersonic flow,
shock waves and expansion waves dominate the flow pattern. Consequently, the pressure distribution over an airfoil in supersonic flow is quite different than that for subsonic flow, giving rise to a major component of drag called wave drag. Supersonic airfoil shapes are characterized by sharp leading edges and thin profiles. Both features reduce the strength of the shock waves and hence reduce the wave drag. See SUPERSONIC FLIGHT. John D. Anderson, Jr. Bibliography. I. H. Abbott and A. E. von Doenhoff, Theory of Wing Sections, 1949; J. D. Anderson, Jr., Aircraft Performance and Design, 1999; J. D. Anderson Jr., Fundamentals of Aerodynamics, 4th ed., 2005; J. D. Anderson, Jr., A History of Aerodynamics, 1997; J. D. Anderson, Jr., Introduction to Flight, 5th ed., 2005.
Airframe The structure consisting of the wings, fuselage, and fin and tail plane (vertical and horizontal stabilizer) of an airplane. This structure has a clearly defined role: to resist all the extreme loads that can be experienced with no danger to the passengers, whether the aircraft be civil or military. The temptation therefore could be to overdesign the airframe and “play safe.” But an overweight civil aircraft would be commercially disastrous, and an overweight military aircraft could be outperformed. To put this matter into perspective, a civil airframe can be up to 40% of the total dry weight, whereas the payload may be as low as 20%. Thus a structure which is overdesigned by only 10% could reduce the payload by 20%. This represents the profit margin on many routes. Design loads. An airframe has to withstand two different types of loading: a maneuver and a sudden gust. A maneuver is intentional, and when a pilot increases the aircraft incidence by applying a download on the tail plane the main wing experiences a sudden increase in lift and accelerates upward. The resulting inertia forces experienced by all the masses can be visualized as balancing the extra lift (Fig. 1). This extra inertia is conventionally referred to as the number of g, where 1 g is the normal weight due to gravity. For an agile combat military aircraft the specification may call for a factor of 9 g or more, depending on the assistance given to the pilot by pressurized suits and whether the maneuver lasts for
lift
masses × ng (n = load factor) Fig. 1. Forces acting on the main wing during a maneuver.
383
384
Airframe many seconds. For a civil aircraft the current allowable acceleration is +2.5 g or −1 g “downwards.” See AERODYNAMIC FORCE. Civil aircraft also have to be designed to resist sudden gusts. A vertical upward gust of velocity u will effectively increase the incidence of the aircraft by u/V radians, where V is the aircraft velocity. Currently the maximum specified gust is 20 m/s (or 45 mi/h). This figure represents the maximum ever recorded in flight. However, all maneuver and gust loadings are multiplied by 1.5 as a safety factor when designing civil aircraft. This somewhat arbitrary figure has withstood the test of time. Structural failures are almost unheard of, provided the specified maintenance is followed. There have been a very small number of cases where pilots have overloaded an aircraft, but nowadays a computer between the pilot and the controls will impose software limitations on control forces. Certain parts of the airframe may be designed to withstand rather special loads, such as the undercarriage in heavy landings. Materials behavior. To discuss airframe design, it is first necessary to describe the behavior of airframe materials. The simplest way to test a material (say, a simple rod) is to pull it and measure the extension due to the applied force. However, a better measure of the actual forces on the material molecules or crystals is the stress (σ ), that is, force divided by crosssectional area. Similarly, a better measure of the crystal deformation is the strain (ε) or fractional change in length of the rod. A final important property is the material’s stiffness (E), which is the ratio of stress to strain for small strains. Some typical stress/strain curves for aircraft materials are shown in Fig. 2. See STRESS AND STRAIN. The most common structural material in use worldwide is mild steel for which the strain is proportional to stress (a linear curve) until “yield,” when the crystals start to move along their boundaries or dislocations. This material then deforms at virtually constant stress, and then when unloaded will do so linearly as shown in Fig. 2, leaving a permanent deformation of, say, 7%. This particular property is ex-
stress (tensile)
carbon composite
light alloy
0.2% proof stress
mild steel
yield stress
0.2
strain (extension), percent
Fig. 2. Stress/strain curves for various types of material.
Fig. 3. Effect of wing box strains of 0.7%.
tremely valuable since it allows thin steel sheets to be pressed into everything from car bodies to domestic cutlery. A similar behavior accurs in aluminum, but by adding small amounts of copper, zinc, magnesium, and so forth, the dislocations can be inhibited at key points. The maximum stress can therefore be increased as shown in Fig. 2. Light aluminum alloys have been the most common airframe material from the early days. To define the maximum allowable stress (in the absence of a true yield), the aircraft sector has concentrated on the allowable strain when the material is unloaded. (Clearly an aircraft should return from flight with the same shape as it took off.) The atypical allowable is a strain of 0.2%, and the corresponding stress is called the 0.2% proof stress. Although these strains may seem very small, an aircraft wing, for example, is not a very efficient structure in bending, and if the strains reached 0.7% the deformations would look something like Fig. 3. See METAL, MECHANICAL PROPERTIES OF; PLASTIC DEFORMATION OF METAL. To obtain higher failure stresses, by avoiding dislocations and extending the linear range, it is necessary to turn to brittle materials like glass, carbon, boron, and various ceramics, all of which have a much purer crystalline structure. However, unlike ductile metals, these materials are very susceptible to flaws and minor cracks. In fact, the standard way that a glazier will cut glass is by inscribing a shallow surface scratch and then propagating it. One way of overcoming this ‘notch sensitivity’ is to manufacture extremely fine fibers and then embed millions of them in a resin matrix. Fibers still have surface flaws, but typically the distance between flaws is about 15– 25 diameters for carbon fibers. Thus, when a single fiber breaks, the load is diffused onto the surrounding fibers since the chances of all fibers having cracks at the same location are zero. A common way of using carbon fiber composites is to make a thin lamina of about 0.13 mm (0.005 in.) thickness in which all the thousands of fibers lie in one direction. A simple laminar (or ply) is extremely stiff and strong in the direction of these fibers but very weak in a direction at right angles to them (like balsa wood). A composite wing skin, for example, is then built by stacking a large number of plies in various directions to suit the nature of the internal stress field. Composite structures are consequently “tailored,” unlike structures of homogeneous metals. Carbon composite structures have dominated combat military aircraft since the 1980s, since they are exceptionally stiff and strong. The table shows
Airframe
385
Properties of some materials used in aircraft structures
Material
E, GPa
σ, MPa
ρ, Mg/m3
E/ρ
σ /ρ
Aluminum alloy Titanium High-tensile steel Carbon composite Glass composite Kevlar composite
71 110 207 200 40 82
480 1000 1720 1500 840 1500
2.8 4.4 7.8 1.55 1.80 1.39
25 25 26 130 22 60
172 226 221 970 460 1080
the values of strength (ultimate tensile strength, σ ) of stiffness (E), and of density (ρ) of some materials used in aircraft structures. The specific stiffness and strength (ratios of stiffness and strength, respectively, to density) are the crucial measures since there is no point in having a superior material that is too heavy. The values in the last two columns of the table indicate that the specific stiffnesses of the three metals are comparable, but clearly carbon composites are superior. Carbon also has better specific strength, matched only by Kevlar. This artificial aramid fiber is cheaper than the other structural materials but is competitive only in tension; it has a very poor compressive strength. See COMPOSITE MATERIAL; MANUFACTURED FIBER. Wing design. An aircraft wing is not a very efficient structure in resisting bending due to gust or maneuver loadings. Clearly, aerodynamic efficiency is the driver, and hence a wing has a small thickness/chord ratio, as low as 1/20 for high-speed military aircraft. When a wing bends due to aerodynamic lift, the top surface goes into compression, the lower into tension. For maximum efficiency, most of the structural material should go to these surfaces, as far as possible away from the center of the section (the neutral axis), where the bending strains are zero. The ubiquitous I section, used in buildings and so forth, is a good example. However, a wing also has to resist torsion, and should twist as little as possible to avoid changing the effective angle of incidence. An I section is hopeless at resisting torsion. For maximum torsional stiffness, a thin-walled structure should be a closed “tube” with as large an enclosed area as possible. In a wing, a fuel tank also fulfills this specification. A typical wing section, with a front and rear spar, is shown in Fig. 4. The trailing edge also contains the flaps and ailerons, and the leading edge may contain deployable slats for increased lift. The upper skin, when in compression, has to be designed to resist buckling. The stress at which a thin plate buckles is proportional to (t/b)2, where t is the plate thickness and b is the width. This property explains the large number of stiffeners shown diagrammatically in
front spar Fig. 4. Typical wing box section.
rear spar
Fig. 5. Wing rib for taking engine mounting on Airbus 380 aircraft. (Courtesy of Airbus UK)
Fig. 4 and intended to reduce the effective width b. The buckling stress is also proportional to (1/l)2, where this time the plate length l is involved and is reduced by inserting ribs into the wing box. These ribs are quite light since their only function is to stop the outer skin’s buckling and to maintain their shape for aerodynamic reasons; a large aircraft may have 40 or 50 of them along the wing. Some ribs are heavily loaded, however; for example, the engine loads have to be diffused into the wing box. Figure 5 shows an engine mounting rib (for the Airbus 380), which is clearly heavy and has an array of stiffeners (machined from the solid) that inhibit buckling under the high shear forces coming from the engine. Where the top edge of the rib is flanged (to be bolted to the wing skin), the cut-outs are seen to allow the wing skin stiffeners to pass through without interruption. See AILERON. The fin (vertical stabilizer) and the tail plane (horizontal stabilizer) have very similar structures to the main wing, although they are much lighter. Actually, their structures should be exceptionally light since they are at the extreme end of the aircraft and if too heavy would adversely move the aircraft’s center of gravity. For this reason all Airbus (and now Boeing) fins and tail planes have been made of carbon composites. A final hazard for the wing designer is the interaction between the wing deflections and the aerodynamic forces: an aeroelastic effect. If the wing lacks sufficient torsional stiffness, it may twist and thereby increase the aerodynamic lift due to the increase in incidence. At a critical divergence speed, the twist may grow beyond limit and lead to structural failure. Another phenomenon, called flutter, occurs when the balance between aerodynamics, structural stiffness, inertia forces, and damping leads to oscillations that may grow without limit. In the absence of aerodynamic forces, the vibration of any structure will gradually subside due to natural material damping. However, the aerodynamic lift and pitching moment are proportional to the surface velocity and result in negative damping. One solution is to introduce mass balancing for controls like the aileron and rudder, so that the extra inertia forces dampen the
386
Airframe oscillations. On a main wing, the forward-placed engines fill this role and ensure a sufficiently high flutter speed. See AEROELASTICITY; AIRCRAFT RUDDER; FLIGHT CONTROLS; FLUTTER (AERONAUTICS); WING; WING STRUCTURE. Fuselage design. The fuselage has a much larger cross section than a wing and is consequently much better able to resist bending and torsion. A fuselage skin may be only 1–2 mm (0.04–0.08 in.) thick, compared to wing skin thicknesses in excess of 20 mm (0.8 in.). Parts of the fuselage also experience bending compression during landing and takeoff, for example. Therefore, as with wings, there is a need for stiffeners and ribs, although the “ribs” are really frames of depth 10–20 mm (0.4–0.8 in.). Heavier frames are needed to take the wing loads into the fuselage, for example, or at undercarriage pickup points. At the rear of the fuselage, the fin or tail plane is often connected to the curved bulkhead, designed to take the cabin pressure loading. Although the working stresses in a fuselage can be low, the design case is really one of fatigue. On every flight, the external pressure at altitude falls to a very low value, while the cabin pressure is maintained at about half an atmosphere for comfort. Thus, in an operating lifetime the fuselage stress field will be cycled thousands of times. All metallic materials can fail at a relatively low stress after a sufficiently large number of cycles. The response of carbon composites is much better since the loading and unloading seems to be reversible with no material damage even at the crystal scale. The presence of windows and doors can exacerbate the fatigue problem by introducing local stress concentrations, but designers now know how to design window and door frames to minimize these concentrations. Unfortunately, this was not the case in the 1950s when the Comet jetliner, the first civil aircraft to cruise at altitudes in excess of 10,000 m (30,000 ft), had fuselage fatigue failures in three aircraft before it was grounded. See FUSELAGE. Composite structures. The superiority of carbon composite materials has already been mentioned, that is, their strength, stiffness, and fatigue performance. In spite of their basic advantages, there are
assumed displacement “shape”
Fig. 6. Finite element approximation.
two disadvantages that have to be overcome, and indeed mostly have been. 1. Carbon composite materials are susceptible to stress concentrations. Laminated thin plates are strong in their own plane but have little resistance to peeling apart of the individual plies. They can easily delaminate at the matrix resin strength. Ideally, all laminates should be in a two-dimensional stress field. But stresses will be three-dimensional near any discontinuity of geometry or load. Careful design and local reinforcements can eliminate this problem. The Airbus 380, for example, has a totally composite center wing box where the two wings intersect the fuselage. This is a region full of joints and discontinuities, and the design is a tribute to the skill of manufacturers and designers who are able to model this problem and analyze the stress fields. 2. Carbon composites can be expensive. Carbon and boron fibers are costly to produce in small numbers. For a long time, the total output of carbon fiber for sporting goods exceeded that needed by the aerospace industry. Only military aircraft could disregard this expense since in that case a poor performer demanded rejection. Not only was the basic material costly, but so were the manufacturing methods. Hand lay-ups of many thin plies (possibly up to 80 in military aircraft) are costly. This cost has gradually been reduced by replacing hand lay-ups with computer-controlled automatic machines, which lay a large number of overlapping tapes with great precision. Another breakthrough has been resin transfer molding, in which a dry fabric can be sandwiched between two faces of a mold and then the resin injected, before cure, under pressure or sucked through under vacuum. The Boeing 787 is an impressive example of a civil aircraft with composite wings and fuselage, having overcome all these problems. Possibly the only threat to the economics of mass-produced airframes is the future legal requirement that manufacturers recycle carbon structures when their operating life is over. Virtual testing. To satisfy either military or civil airworthiness authorities, aircraft have to be tested to ultimate design loads. Prior to this, the manufacturer will have made many component tests on panels, ribs, undercarriages, and so forth. These very expensive and time-consuming tests have been traditional since no theories were available for highly three-dimensional and detailed components. However, since the 1980s remarkable progress has been made in the ability to simulate in a computer model the behavior of structures loaded to failure. Basically there are only three sorts of equations to satisfy in theoretical structures: (1) the stresses everywhere have to be in equilibrium with themselves and the applied forces; (2) the strains have to be related to the displacement field (purely a geometrical argument); and (3) the two are linked by a stress/strain law based on material tests. The main problem has been the enormous geometrical complexity of aircraft structures at the detailed level, but the advent of the digital computer (and later
Airglow
X
Z
Y Fig. 7. Finite element model of a wing box. (Courtesy of Airbus UK)
computer-aided-design software) has changed this forever. The finite element method is conceptually simple. A typical displacement field across a wing from leading to trailing edge is illustrated in Fig. 6. The structure is then divided up into finite elements which are so small that a reasonable approximation can be made for the shape of the variation across any element. In Fig. 6 a linear variation is implied. Then, all that remains is the magnitude of the displacements at, say, the edges of the elements. The strains can then be expressed in terms of these magnitudes, the strain/stress law implemented, and finally the equilibrium equations. The result is a large number of equations giving the many displacements in terms of the applied forces. The number of displacement unknowns and equations can run into hundreds of thousands, but modern computers are powerful and inexpensive enough to cope with problems of this magnitude. A typical finite element display is shown in Fig. 7. See COMPUTER-AIDED DESIGN AND MANUFACTURING; FINITE ELEMENT METHOD. Commercial computer codes have been available since the 1970s, but only recently has it been feasible to simulate the failure process, which is nonlinear as damage propagates and which modern codes can simulate. Virtual testing is now a reality and more informative than field testing since the internal stress fields are modeled in detail—for example, when a bird strikes an aircraft or is sucked into an engine fan. See AIRCRAFT TESTING. Glyn A. O. Davies Bibliography. H. D. Curtis, Fundamentals of Aircraft Structural Analysis, Times-Mirror HE Group, Los Angeles, 1997; M. Davies (ed.), The Standard Handbook for Aeronautical and Astronautical Engineers, McGraw-Hill, New York, 2003; T. H. Megson, Aircraft Structures for Engineering Students, 3d ed., Edward Arnold, London, 1999.
Airglow Visible, infrared, and ultraviolet emissions from the atoms and molecules in the atmosphere above 30 km (20 mi), generally in layers, and mostly between 70 and 300 km (45 and 200 mi). The airglow, together with the ionosphere, is found in the uppermost parts of the atmosphere that absorb the incoming energetic radiations from the Sun. While the airglow consists of spectral features similar to those of the aurora, it is mostly uniform over the sky; and it is caused by the absorption of solar ultraviolet and x-radiations, rather than energetic particles. See AURORA. Dayglow. The daytime airglow (dayglow) is caused mainly by fluorescence processes as molecules and atoms are photodissociated and photoionized. The photoelectrons that are produced in the ionization processes are a further source of airglow in their collisions with other atoms and molecules. Emissions produced by recombination of previously produced neutral species and ions are a minor source of dayglow, and resonant and fluorescent scattering of sunlight also contributes. See FLUORESCENCE. Only during twilight is the blue sky glow that is caused by sunlight scattered on the lower atmosphere sufficiently weak that some dayglow emissions can be observed from the ground. Thus, twilight offers an opportunity to observe resonant scattering of sunlight on layers such as those of the alkali atoms sodium, lithium, and potassium. As the Earth’s shadow scans through the layers, the changes of intensity allow their heights (near 90 km or 55 mi) to be measured. The atoms can be observed by lidar (laser radar) with similar results but more detail. See ALKALI EMISSIONS; LIDAR; SCATTERING OF ELECTROMAGNETIC RADIATION.
387
388
Airplane Nightglow. The nighttime airglow (nightglow) is predominantly due to recombination emissions. The ionospheric plasma recombines near the bottom of the F region (150–200 km or 90–120 mi) where the densities and thus collision frequencies are higher, producing bright atomic oxygen (O) spectral lines in the red (at 630 and 636 nanometers) and a weaker green (558-nm) line. The recombination of the neutral radicals that are generated in daytime photodissociation proceeds fastest in the mesosphere, where the densities are higher at the bottom of the region of production. Atomic oxygen recombination forming excited dioxygen (O2) molecules leads to several characteristic spectral-band systems. The recombination energy can alternatively be transferred to an oxygen atom in collision and can become the source of a strong green (558-nm) emission. Excited oxygen atoms that could produce the red (630- and 636-nm) lines in the mesosphere are nearly always rapidly quenched in another collision. See ATOMIC STRUCTURE AND SPECTRA. Strong infrared hydroxyl (OH) emission arises in the mesosphere from the reaction of atomic hydrogen (H) with ozone (O3), and the H is regenerated by the OH reacting with O to produce O2, with the net result being the recombination of O and O3 and the recycling of the H. Similar cyclic reactions excite sodium atoms. A weak continuum in the green is due to the recombination of nitric oxide (NO). The total of all the visible nightglow emissions, together with zodiacal light and scattered starlight, can be seen with the naked eye as the faint light between stars. From space, the edge-on view of airglow is of a band of light above the horizon. See ATMOSPHERIC OZONE; ZODIACAL LIGHT. Observations. Observations made by spacecraft are free of the constraint of the blue sky glow, and can also detect the ultraviolet emissions that cannot be seen from the ground. Such observations of resonant ultraviolet scattering of oxygen ions (O+) and helium ions (He+) can be used to map the distribution of ions in the plasmasphere to distances to four times the Earth radius. Scattering in the Lyman series maps the so-called geocorona of atomic hydrogen out to 15 times the Earth radius. A variety of remote-sensing techniques from the ground and from space utilize airglow to determine the composition and structure not only of the Earth’s upper atmosphere but also of the atmospheres of other planets. Spectroscopic observations of relative intensities and Doppler profiles identify the composition, temperature, winds, and energy inputs, as well as ionospheric layer heights and ion concentrations. Images of extended areas of terrestrial airglow show wave structures caused by atmospheric gravity waves. These waves originate deep in planetary atmospheres, and their amplitude grows exponentially with height; their dissipation is an important source of heating of upper atmospheres. Images made of the ionospheric recombination emissions, looking toward the magnetic equator from low-latitudes sites on Earth, show large-scale structures due to plasma instabilities.
The airglows from other planets are remarkably different from that of Earth. The airglows from Mars and Venus show strong emissions of carbon-oxygen species (CO, CO+, CO2+), O, and H; from Titan emissions of molecular nitrogen (N2) and atomic nitrogen (N); and from Jupiter and Saturn, emissions of molecular hydrogen (H2) and H. These emissions confirm that CO2, N2, and H2 are the dominant constituents in each case. Brian A. Tinsley Bibliography. J. W. Chamberlain, Physics of the Aurora and Airglow, 1961; J. W. Chamberlain and D. M. Hunten, Theory of Planetary Atmospheres, 2d ed., 1987; G. Paschmann et al., Auroral Plasma Physics, 2003; M. H. Rees et al., Physics and Chemistry of the Upper Atmosphere, 1989.
Airplane A heavier-than-air vehicle designed to use the pressures created by its motion through the air to lift and transport useful loads. Although airplanes exist in many forms adapted for diverse purposes, they all employ power to overcome the aerodynamic resistance, termed drag, thereby achieving forward motion through the air. The air flowing over specially designed wing surfaces produces pressure patterns which are dependent upon the shape of the surface, angle at which the air approaches the wing, physical properties of the air, and velocity. These pressure patterns acting over the wing surface produce the lift force necessary for flight. See AIRFOIL. To achieve practical, controllable flight, an airplane must consist of a source of thrust for propulsion, a geometric arrangement to produce lift, and a control system capable of maneuvering the vehicle within prescribed limits. Further, to be satisfactory, the vehicle should display stable characteristics, so that if it is disturbed from an equilibrium condition, forces and moments are created which return it to its original condition without necessitating corrective action on the part of the pilot. Efficient design will minimize the aerodynamic drag, thereby reducing the propulsive thrust required for a given flight condition, and will maximize the lifting capability per pound of airframe and engine weight, thereby increasing the useful, or transportable, load. Air propulsion systems. Because the airplane moves through a fluid medium, its propulsion system must produce a change in the momentum of its working fluid to generate the thrust force required to overcome aerodynamic drag. Reciprocating and turboprop engines produce power by converting chemical energy of fuel into rotation of a shaft; a propeller fitted to this shaft accelerates the air passing through its rotational area, thereby changing engine power into useful thrust. Turbojet and ramjet engines produce thrust directly by adding heat to the air passing through the engines and ejecting the resulting hot exhaust gases at high velocity. A fanjet (turbofan) engine is a hybrid between the turboprop and the pure turbojet that combines a number of features common to both. The turbojet, ramjet, and fanjet en-
Airplane gines use the surrounding air both as primary working fluid and as oxidizer to support combustion. See JET PROPULSION; PROPELLER (AIRCRAFT); RAMJET; RECIPROCATING AIRCRAFT ENGINE; TURBOFAN; TURBOJET; TURBOPROP. Thrust. The engine-propeller combinations, whether reciprocating or turboprop, are effective at low forward speeds, but their thrust falls off rapidly as flight speeds approach the speed of sound (Mach 1), largely because of the tremendous increase in power required to turn the blades when shock waves form at their tips. The turbojet thrust increases with speed because of aerodynamic compression, which becomes more marked at higher speeds. The fanjet, by combining many of the characteristics of the propeller and the turbojet, provides efficient propulsion in the range of speeds between those where the thrust of the propeller falls off rapidly owing to the compressibility drag rise at the tips, and the thrust of the turbojet picks up. The limit of turbojet operation is set largely by the temperature limitations of the materials used to fabricate the turbine and by shock patterns which may produce unsteady flow and separations at high speeds, interfering with efficient aerodynamic compression, creating difficulties with the intake system, and causing the engine to stall and flame out. See SUPERSONIC DIFFUSER. The ramjet cannot produce thrust at low forward speeds because it depends upon an appreciable aerodynamic (ram) compression. It therefore requires a booster system to raise the airplane speed sufficiently that useful thrust can be produced. Because of improved aerodynamic compression, as with the turbojet, the thrust produced increases as the velocity of flight increases. Since a ramjet is designed for maximum performance at a given flight condition, its intake design may be such that its range of operating speeds is small. Weight and fuel consumption. Ramjet engines, although light, have large fuel consumptions, whereas the heavy reciprocating engine has the lowest consumption, and turboprop, fanjet, and turbojet engines are intermediate in weight and fuel consumption (see table). Requirements of speed, range, size, and cost must be balanced to select the engine best suited for the specific airplane design. See AIRCRAFT ENGINE PERFORMANCE; SPECIFIC FUEL CONSUMPTION. Aerodynamic resistance. Aerodynamic resistance is composed of four parts: (1) skin friction drag, which
arises from the viscosity of the air; (2) pressure drag, which arises from the pressure field about the body; (3) wave drag, which arises from the compressible nature of air, and which occurs as flight speeds approach and exceed the speed of sound; and (4) induced drag, which is a resistance produced by the generation of lift. Skin friction. The magnitude of skin friction drag depends upon the surface area of the airplane and the nature of the flow layer that is in immediate contact with it. If the boundary-layer flow is smooth and steady (laminar), the skin friction generated is much less than when this layer becomes turbulent. As a laminar boundary layer proceeds along a surface, depending upon the smoothness and shape of the surface, as well as the density, viscosity, and velocity of the air, the fluid momentum of the flow decreases. Simultaneously, disturbances in the flow or the surface produce motions perpendicular to the surface; these motions tend to grow as the fluid momentum decreases. Eventually a transition condition is reached at which these motions cause the laminar boundary layer to change to turbulent flow. Hence, unless special devices are employed, such as suction slots to remove the low momentum boundarylayer air before it becomes turbulent, the laminar flow will break down after it has passed over a sufficient length of surface. For this reason, on those portions of the airplane where flow initiates, such as the nose and wing leading edges, extreme smoothness to minimize viscous losses is desirable, while further aft, where the flow is turbulent, greater surface roughness can be tolerated. See BOUNDARY-LAYER FLOW. Pressure drag. If sufficient momentum is lost in the boundary layer, the main flow breaks away from the surface, forming a region of separation. Separation results in considerably lower surface pressures than would result in the ideal streamline case. As a consequence, suction pressure on the downstream portions of the body produces pressure drag. This drag is minimized by shaping the body geometrically to reduce the separation region to a minimum by giving it a streamlined shape (Fig. 1). In modern high-speed aircraft, pressure drag has been reduced to a small percentage of the total drag. See STREAMLINING. Wave drag. As flight speed of the airplane increases, there comes a point where the flow passing over
Fuel consumption and weight of principal types of airplane engine
Engine
Fuel consumption, lbm fuel/(h)(lbf thrust)∗
Weight, lbm engine/ lbf thrust†
Ramjet Turbojet Fanjet (turbofan) Turboprop Reciprocating
1.7–2.6 0.75–1.1‡ 0.3–0.5 0.4 0.2
0.4 0.8§ 0.7§ 1.2 2.4
* Minimum values. Value for ramjet is for Mach 2; values for turbojet and fanjet engines are based on static ground tests; values for reciprocating and turboprop engines are for speeds below Mach 0.3. 1 lbm/(h)(lbf) = 0.1 kg/(h)(N). † Based on total installed nacelle and power-plant weight, and on thurst at 375 mi/h (168 m/s) at sea level. 1 lbm/lbf = 0.1 kg/N. ‡ Value for Olympus 593 engine aboard Concorde is 1.19 at Mach 2. § Approximate cruise value.
389
Airplane
region of separation streamlines
turbulent wake
(b) Fig. 1. Effect of streamlining. (a) Body without streamlining. (b) Body with streamlining, minimizing the region of separation.
a surface reaches the local speed of sound. At this speed (Mach 1) pressure waves propagate through the air. This wave propagation forms shock patterns that interact with the boundary layer to cause premature separation and a sudden drag rise. Because, to create low pressures necessary to generate lift, local velocities next to the upper surface of the wing can be appreciably higher than the flight speed, this separation usually takes place over the wing. Separation over a wing causes an abrupt increase in drag and loss of lift. As the speed continues to increase, shock patterns become steady, the bow and stern shocks changing location and becoming much stronger until they become analogous to the surface waves that are created by a ship. As with such waves, at supersonic speeds the generation of these shock waves represents the major contribution to the drag of the vehicle. See AERODYNAMIC WAVE DRAG. Early experimenters discovered that it was possible to delay and alleviate somewhat the drag rise associated with the shock-induced separations that occur when local sonic speed is first reached at some point on the wing surface. They swept back the leading edge of the wing with respect to the flow direction. The velocity past the wing can be considered as comprising a component in the spanwise direction and a component normal to the span (Fig. 2). Only this latter component affects lift and is subjected to local accelerations produced by the airfoil shape. As a consequence, the free-stream velocity can be considerably increased over the value that would produce force divergence for a straight wing. When the flight Mach number becomes sufficiently large, the effect of sweep tends to vanish. Other attempts to delay the onset and lessen the magnitude of shock-induced separations have concentrated on airfoil design. If the region in which the air passing over the upper surface exceeds the
true wind velocity (V )
stagnation point
Λ
(a)
os
vortices
speed of sound can be distributed over more of the surface at a given flight speed, the magnitude of both the shock and the resulting separation can be reduced. So-called supercritical airfoils offer the possibility either of increasing the Mach number (speed of flight/speed of sound) by the order of 15% or, holding the Mach number constant, of increasing the thickness of the wing by as much as 40%. See SUPERCRITICAL WING. Area rule. Studies, at first theoretical in nature but later confirmed by experiment, demonstrated that the drag rise associated with the initial formation and subsequent growth of the shock wave pattern about an aircraft traveling at transonic and supersonic speeds was a function of the rate of change of aircraft cross-sectional area. Furthermore, these studies showed that, if there were an abrupt change in the cross-sectional area perpendicular to the line of flight, the drag rise at a given speed would be greater than if the change were tailored to be more gradual. For this reason aircraft are now designed with the fuselage shape, as well as engine nacelle location and shape, so arranged to produce optimum variation of cross-sectional area for the given design speed; the design principle is termed the area rule. (Fig. 3). See TRANSONIC FLIGHT. Induced drag. The final major type of aerodynamic resistance that must be considered arises from the fact that in creating lift the wing must impart a downward momentum to the air flowing past. This downward component of local velocity, termed downwash, rotates the resultant force of the wing in such a manner that a force component is directed in the downstream direction. This component, which depends directly upon the magnitude of the lift force, is termed the induced drag. Lift of an airplane wing depends upon its shape, speed of flight, angle the wing makes with the free
Vc
stagnation point
spanwise flow
root chord
390
leading edge
angle of sweepback (Λ) quarter chord
trailing edge tip chord Fig. 2. Geometry of an airplane wing. Wind velocity component perpendicular to the wing span (V cos Λ, where V is the true wind velocity and Λ is the angle of sweepback) produces lift.
Airplane 1.6
L = 1/2 ρV 2 SCL
1.2
(1)
V are the density and speed of the air, and S is the wing planform area. The coefficient of lift is generally
coefficient of lift (CL)
stream, and properties of the air. These factors are all grouped in the nondimensional coefficient of lift CL, which is related to the lift L by Eq. (1), where ρ and
0.8
0.4 A
A _8
_4
0
4
8
12
16
angle of attack (α) B
B
Fig. 4. Variation of coefficient of lift with the angle of attack, showing increase up to the stall angle.
C
C
uncompensated body
area-compensated body
(a)
section A
section A
section B
section B
section C uncompensated body
section C area-compensated body
(b)
cross-sectional area
A
B
C uncompensated body
area-compensated body body station
coefficient of drag (CD )
(c)
uncompensated body
area-compensated body Mach number
(d) Fig. 3. Illustration of the area rule. (a) Half-plan views of airplanes with uncompensated body and areacompensated body. (b) Cross sections of these airplanes. (c) Variation of cross-sectional areas of these airplanes with body station (longitudinal distance along aircraft). (d) Resulting variation of drag coefficients CD of these airplanes with Mach number.
presented as a function of the angle of attack (Fig. 4). To maintain level flight, the wing must produce a lift force equal to the weight of the airplane. For a given altitude this implies a low CL at high speeds but, as the speed is decreased, CL must increase if level flight is to be maintained. The pilot actually controls the speed of the airplane by changing the angle of attack. There is a maximum angle of attack beyond which no lift increase is experienced (Fig. 5). As this angle is approached, the air in the boundary layer next to the wing surface experiences such energy losses that it can no longer flow over the wing, but separates, resulting in a loss of suction pressures and a decrease in lift. When there is no further increase in lift as the angle of attack increases, the wing is said to be stalled. Because induced drag depends upon the magnitude of the lift coefficient, it is of concern primarily in those flight regimes where a large lift coefficient is encountered: relatively low speed or high altitudes. The magnitude of the angle through which the resultant force of the wing is rotated is a function of the aspect ratio, the ratio of the span to the chord of the wing. The higher this ratio, the smaller the angle. Hence airplanes designed for flight at high altitude or for cruising at relatively low speeds will have long, slender wings. Planes designed for higher speeds will be less affected by the induced drag and will generally have much smaller spans to minimize the other forms of resistance. A measure of efficiency of the airframe is the ratio of lift to drag that it can produce. To achieve very low speeds, special devices such as flaps or slots are used. Such devices, as well as high-lift boundary-layer control, where separation of the boundary layer is delayed either by sucking the low-momentum air next to the surface away through slots or holes in the surface or by adding momentum to the layer by blowing into it through slots, increase the lift that can be achieved before the wing stalls. These devices are thus useful for landing and takeoff maneuvers. The increase of lift achieved by their deployment is
391
392
Airplane usually accompanied by an appreciable increase of drag and hence does not increase airframe efficiency. See AERODYNAMIC FORCE; AERODYNAMICS. Range. The distance that an airplane can travel with a given amount of fuel is its range, expressed by Eq. (2), where Ct is specific fuel consumption in V L dw Range = (2) Ct D w pounds of fuel consumed per pound of thrust per hour, V is flight velocity, L/D is lift to drag ratio, and dw/w is a measure of the rate at which total airplane weight changes during the flight. If the total difference in weight during the flight is produced by fuel consumption, the term (V/Ct)(L/D) should be maximized. In general, it is desirable to have low specific fuel consumption, high speed, and high L/D. Unfortunately, these factors do not maximize at the same point. With reciprocating enginepropeller combination aircraft, maximum range could be achieved at the speed for maximum L/D. With the advent of other forms of propulsion, however, the speed for maximum range has increased. With some aircraft it is possible to achieve high range even at supersonic speeds, for although L/D is reduced, the increase in flight velocity tends to compensate. Airplanes traveling at supersonic speeds create shock waves which extend outward like the wake of a ship. As the speed and size of aircraft increase, the impact of the shock waves becomes increasingly noticeable and possibly destructive in the form of a sonic boom. Commercial airplanes designed to fly at supersonic speeds must, therefore, operate to reduce shock waves when flying over populated areas. The airplanes must fly either at extremely high altitudes or at speeds below the speed of sound. A significant amount of flight time of supersonic aircraft is spent flying at speeds below Mach 1. The L/D of typical optimum supersonic configurations with highly swept wings and low aspect ratio is low at subsonic speeds. The result of the combination of low L/D and low speed is a substantial reduction in range. To overcome this disadvantage, airplane configurations capable of changing the sweep of the wings in flight have been developed. The variable-geometry compromise provides both high-aspect-ratio, nearly straight wings at low speed and low-aspect-ratio, highly swept wings at supersonic speeds. See SONIC BOOM; SUPERSONIC FLIGHT. Stability. A necessary element of a successful airplane is stability, which is the tendency of forces and moments to be set up that restore the vehicle to its equilibrium position if it is disturbed from this position. For an airplane to be in equilibrium, the sum of the forces acting at its center of gravity must equal zero, as must the sum of the moments. See FLIGHT CHARACTERISTICS. To examine the condition of stability, suppose that the equilibrium of the airplane is suddenly disturbed in a manner that causes it to accelerate. This means that its lift coefficient decreases (Fig. 5). A stable airplane would immediately be subjected to a nose-up
moment which would tend to increase the lift coefficient and reduce the speed. Similarly, a reduction in the speed of a stable airplane should produce a nose-down moment. Stability is achieved by balancing the force and moment contributions of various portions of the airplane so that the summation is stable. In conventional configurations the major stabilizing contribution arises from the tail, but by careful design tailless or tail-first (Canard) configurations are possible. Because the balance is achieved for moments about the center of gravity, the location of this point is important. Its position depends upon the load condition and can change in flight as fuel is consumed or as disposable load is dropped. The airplane must be designed so that it reacts in a stable manner throughout the entire range of center-of-gravity locations encountered in service. Controls. An airplane must be controllable. Major aerodynamic control surfaces include the elevators, ailerons, and rudder (Fig. 5a). The angle of attack, and consequently the speed in level flight, is controlled by the elevators. Deflecting these surfaces downward increases the lift of the horizontal tail (just as flaps increase the lift of the wing) and hence causes the airplane to nose downward. An upward deflection produces the opposite effect. See ELEVATOR (AIRCRAFT). Lateral control is achieved by differential action of ailerons. One aileron goes down, increasing wing lift on that side, and the other goes up, reducing lift on that side. The effects couple to produce a rolling moment that tilts the wing. Ailerons also serve to turn the airplane by tilting the wing lift and producing a component normal to the flight path tending to curve it (Fig. 5b). See AILERON. Occasionally ailerons are replaced by a spoiler, a device which projects upward on the downgoing
stick
rudder
rudder pedal
elevator (a)
aileron L
(b) Fig. 5. Aerodynamic control of an airplane. (a) Basic controls and control surfaces. (b) Effect of tilting the airplane in the direction of the wing lift L.
Airplane wing, thus producing separation that reduces lift. The spoiler may be used for several reasons. It reduces the twisting moment, which for an aileron can be so great that the torsional resistance of the wing is exceeded, causing it to rotate and thereby changing its local angle of attack. The upgoing aileron, instead of reducing lift, can produce such a change in angle of attack that the control action is reversed. A spoiler is also used to overcome adverse yaw, the tendency of the higher lift on the upgoing wing to produce an increased induced drag on that side. This drag tends to turn the airplane in the wrong direction. The vertical tail provides directional stability. The rudder prevents the airplane from slewing, or picking up too large an angle of yaw during maneuvers. Although the airplane could be turned with the rudder, such a maneuver would be uncomfortable, and the ailerons are always used. If the vertical tail were large enough to prevent slewing by itself, the rudder could be eliminated and a two-control machine built. See AIRCRAFT RUDDER; FLIGHT CONTROLS. Augmentation systems. While it is both possible and desirable to build aircraft that are inherently stable at relatively low flight speeds, as these speeds increase into the transonic and supersonic ranges the magnitudes and nature of the pressure fields about the aircraft change, resulting in two effects. The first arises from the airloads becoming so great that not only do the forces required to move the controls exceed the pilot’s strength but the aircraft, normally thought of as a rigid body, distorts and twists, changing its aerodynamic characteristics (aeroelastic effects). The second effect results from the change in the nature of the distribution of the loads over the airframe and control surfaces (compressibility effects). In combination the two effects alter both the stability and controllability characteristics of the airplane. See AEROELASTICITY. To handle such cases it is often necessary to provide systems that not only supply the necessary power to move the control surfaces but also augment the stability characteristics of the basic aircraft by sensing motions and automatically applying compensating control inputs independent of pilot action. A common application of such a system is to augment yaw damping at high subsonic speeds. In many modern designs, particularly military ones, the pilot is no longer physically connected to the control surfaces. Instead, pilot control motions are interpreted by a computer that superimposes pilot commands on those control motions required to maintain stability. In some cases the airplane has been designed to be unstable if unaugmented in order to achieve high maneuverability. See STABILITY AUGMENTATION. Structures. Airplane structures must provide the required strength and rigidity, within the shape envelope dictated by the aerodynamic requirements of the airplane, for the lowest possible weight. Other factors such as manufacturing ease, corrosion resistance, fatigue life, and cost are important in selecting both the materials and the nature of the structural system to be employed, but the strength-to-weight ratio is of critical importance.
Most materials used in aircraft construction display the characteristic that stress is essentially proportional to strain up to a point called the proportional limit. Removing the load before this point is reached returns the material to its original unstrained shape, while exceeding the proportional limit results in permanent distortion. Airplane structures are designed so that no stress exceeds the proportional limit if the airplane is maneuvered over its full range of speed and design accelerations: the ultimate stress is not exceeded even if the structure is subject to 11/2 times its design load. See AIRFRAME; FUSELAGE. The major structural problem is design of the wing. Aerodynamic considerations demand that wing thickness be kept small with respect to wing chord. The air load is distributed along the wing, and the magnitude and configuration of the load are a function of the maneuver under consideration. The wing may thus be considered as a long, slender beam subjected to loads that during accelerations may total several times the weight of the aircraft. Because the depth of the beam (airfoil thickness) is small, the outer fiber stresses are large. See WING; WING STRUCTURE. One solution to the problem of high bending stresses is the biplane, where the upper and lower wings, which are connected by struts and wires, act as the flanges of a beam having a depth equal to the separation between the wings. Another solution is the use of a pin connection at the wing root and an external strut. High-performance airplanes cannot afford the drag penalties of such solutions and must employ cantilever wings. Because of the high bending and torsional stresses encountered, a cantilever structure implies a stressed skin, or outer covering of the wing. The skin may be thin, supported by many stringers, or relatively thick with fewer longitudinal members. Ribs along the span aid in distributing loads and in reducing column action in the stringers. High flight speeds demand smooth skin surfaces. These are sometimes produced by using skin milled from metal blocks rather than sheet. Two thin sheets bonded to a low-density core form a sandwich, a type of construction frequently used. Aircraft flying at high speeds encounter high temperatures due to the high aerodynamic compression that occurs at the nose and the leading edges of the wings. At high Mach numbers these temperatures become sufficiently high to reduce the strength of aluminum and its alloys, so other materials, such as stainless steel or titanium, are sometimes used in these areas. When flight speeds become so high that the temperature rise affects the characteristics of the construction material, it may be necessary to employ cooling in sensitive areas. See AEROTHERMODYNAMICS. David C. Hazen Bibliography. J. D. Anderson, Jr., Introduction to Flight, 5th ed., 2005; R. Jackson (ed.), Jane’s All the World’s Aircraft, revised periodically; M. J. H. Taylor and D. Mondey, Guinness Book of Aircraft, 6th ed., 1992.
393
394
Airport engineering
Airport engineering The planning, design, construction, and operation and maintenance of facilities providing for the landing and takeoff, loading and unloading, servicing, maintenance, and storage of aircraft. Facilities at airports are generally described as either airside, which commences at the secured boundary between terminal and apron and extends to the runway and to facilities beyond, such as navigational or remote air-trafficcontrol emplacements; or landside, which includes the terminal, cargo-processing, and land-vehicle approach facilities. The design of airports involves many diverse areas of engineering. Their significant cost has resulted in the development of sophisticated techniques for planning, construction, and operation of airports. Design considerations. Airport design provides for convenient passenger access, efficient aircraft operations, and conveyance of cargo and support materials. Airports provide facilities for changing transportation modes, such as people transferring from cars and buses to aircraft, cargo transferring from shipping containers to trucks, or regional aircraft supplying passengers and cargo for intercontinental aircraft. In the United States, engineers utilize standards from the Federal Aviation Administration (FAA), Transportation Security Administration (TSA), aircraft performance characteristics, cost benefit analysis, and established building codes to prepare detailed layouts of the essential airport elements. These elements include airport site boundaries, runway layout, terminal-building configuration, supportbuilding locations, roadway and rail access, and supporting utility layouts. Airport engineers constantly evaluate new mechanical and computer technologies that might increase throughput of baggage, cargo, and passengers. Design choices must enhance user orientation and comfort. Because airports and airlines are key transportation and economic facilities for major cities worldwide, airport engineers frequently integrate a wide variety of secondary design considerations, including air- and water-quality enhancement, noise concerns, citizen participation, wildlife impacts, and esthetic features. Site selection for new or expanded airports. Site selection factors vary somewhat according to whether (1) an entirely new airport is being constructed or (2) an existing facility is being expanded. Few metropolitan areas have relatively undeveloped acreage within reasonable proximity to the population center to permit development of new airports. For those airports requiring major additional airfield capacity, and hence an entirely new site, the following factors must be evaluated for each alternative site: proximity to existing highways and major utilities; demolition requirements; contamination of air, land, and water; air-traffic constraints such as nearby smaller airports facilities; nearby mountains; numbers of households affected by relocation and noise; political jurisdiction; potential lost mineral or agricultural production; and cost associated with all these fac-
tors. Some governments have elected to create sites for new airports using ocean fills. The exact configuration of the artificial island sites is critical due to the high foundation costs, both for the airport proper and for the required connecting roadway and rail bridges. New airports offer significantly increased runway capacity over airport reconstruction projects due to the larger space available for more runways, longer runways, and more widely spaced runways. New airports are rare: Dallas/Fort Worth International Airport opened in 1972, and Denver International opened in 1995. In 1998, two new airports constructed at Air Force bases opened: the Austin, Texas, airport; and MidAmerica, at Scott Air Force Base in Illinois. Killeen Airport in Texas opened as a jointuse facility with Robert Gray Army airfield in 2004. Denver capitalized on its geographic location in the midsection of continental America, its historic hub airline operations, and undeveloped land adjacent to the metropolitan area to build an unusually large airport (Fig. 1). This location is farther from the Rocky Mountains than the old airport, a significant aeronautical advantage. The number of citizens affected by the FAA’s regulated noise level was significantly reduced by the new airport’s location. Denver’s master plan and environmental approvals were completed between 1986 and 1989. Design, construction, training, testing, and airport certification were completed between 1989 and 1995. For commercial airport sites being expanded or military bases being revamped for commercial use, many site factors must be considered. The possibility of acquiring adjacent commercial or residential properties requires extensive analysis of social and economic impacts and development costs. Mitigation measures of noise impact constitute a key element of any proposed airport expansion. Such measures can include restrictions on types of aircraft using particular runways, avoidance of noise-sensitive areas by air- traffic routing, acquisition of additional land adjacent to the airport as a noise buffer, improvements to existing occupied structures to dampen noise transmission, or construction of noise berms for ground operations. Airfield configuration. Since the runways and taxiways constitute the largest portion of the airport’s land mass, their layout is generally one of the first steps in the airport planning process. The optimum runway configuration, or layout, is based on longrange forecasts of the estimated numbers of aircraft landings and departures, airport elevation, local prevailing wind direction, and regional and national airtraffic patterns. A paved runway surface 12,000 ft (3660 m) long and 150 ft (45 m) wide is suitable for most applications. However, runway length requirements vary according to the type of aircraft, temperature, and elevation. A parallel taxiway is generally constructed 600 ft (180 m) from the runway (measured centerline to centerline). It is connected by shorter high-speed taxiways to allow arriving aircraft to leave the runway surface quickly so that another aircraft can land. This combination of paved
Airport engineering 0
1
5 mi
2
NORTH
TO BOULDER
LAFAYETTE
BRIGHTON PR
LOUISVILLE
OP
OS
ED
E-
47
0
25 85
THORNTON
BROOMFIELD
76
WESTMINSTER NORTHGLENN
FEDERAL HEIGHTS
36
WILDLIFE PRESERVE 76
ARVADA 270
76 70
FOOTHILLS
GOLDEN
PENA BLVD.
COMMERCE CITY STAPLETON INTERNATIONAL AIRPORT
70
WHEAT RIDGE EDGEWATER
DENVER INTERNATIONAL AIRPORT
DENVER 70
COLFAX AVE.
DOWNTOWN DENVER
6TH AVE.
LAKEWOOD
225
AURORA
GLENDALE BROADWAY
25
285
BOW MAR
SHERIDAN ENGLEWOOD LITTLETON
470
CHERRY HILLS VILLAGE
LINE OF BUILTUP AREA
GREENWOOD VILLAGE
HIGHLANDS RANCH
Fig. 1. Regional map of the Denver International Airport.
surfaces is generally referred to as a runway-taxiway complex. Airports provide runway layouts so that incoming and departing aircraft can maintain required minimum separation for safe, simultaneous aircraft operations. For parallel runways, thresholds would be slightly staggered to avoid wake turbulence interference between incoming aircraft. Staggered thresholds might also be used to minimize crossing of active runways by taxiing aircraft. Each runway crossing is a potential aircraft delay and a safety hazard. When airports have sufficiently high-velocity crosswinds or tailwinds from more than one direction, crosswind runways are located at some angle to the primary runway as dictated by a wind rose analysis. See WIND ROSE. Alternative airfield configurations are generally analyzed under various operating conditions by using computer simulations. Typically, these simulations assume a certain number of aircraft arrivals and departures, based on actual airline schedules for an entire year, weather conditions based on historical data, and various combinations of runway usage by the air-traffic controllers; and then the required time for each aircraft operation is calculated (Fig. 2). The configuration with the least amount of arrival time and taxi time on an annual basis is preferred. Justification for additional runways would accrue as the construction and financing costs of the runway are
compared to the operational savings of reduced aircraft delays. Some runways, widely spaced so that the aircraft operations are independent, are used primarily for instrument flight rules (IFR) conditions, when visibility is low. In visible flights rules (VFR) conditions, FAA rules allow independent operations with more closely spaced runways. The result is shorter taxi distances after the aircraft lands. The IFR runways, which are farther from the terminal, might also be used during peak periods to increase the airport’s capacity. If excess runway capacity is available, FAA air-traffic controllers have more discretion to assign aircraft landings and departures to minimize noise exposure and taxi distances. In the United States, the FAA defines standard techniques for calculating the hourly capacity, annual service volumes, delays, average taxiing distances, and annual runway crossings. After the runway configuration is determined, detailed construction drawings are prepared that dictate the runway profile and cross section, pavement depth and materials, subsurface remediation, drainage controls, lighting, and paint markings. A cost-effective design balances grading materials from the airport site to minimize trucking costs and the associated impact on air quality and traffic congestion, and specifies locally available materials for the paving mix. See AIR POLLUTION; PAVEMENT.
395
396
Airport engineering
Airfield Layout No. 1 Operating Configuration No. 1
Airfield Layout No. 1 Operating Configuration No. 2
Airfield Layout No. 2 Operating Configuration No. 1
Airfield Layout No. 2 Operating Configuration No. 2 Key: arrival
departure
Fig. 2. Combinations of arrivals and departures on the same runway configuration.
Runways are paved with concrete, asphalt, concrete-treated base, or some combination of layers of these materials (Fig. 3). Runways for larger aircraft require thicker, more expensive pavement sections. Engineers design these pavements for longterm durability. The expected life of a concrete runway can be increased from 20 to 40 years, based on enhanced mix designs and sections. See CONCRETE. The longitudinal profile of a runway requires an iterative design process. The objective is to minimize grade changes for the smoothest aircraft operations possible, while minimizing the total quantities of earthwork (cut-and-fill) to construct the runway. The availability of computerized controls and improved lighting fixtures is increasing the use of technology known as surface movement guidance systems for pilots. These advanced airfield lighting systems provide positive control for taxiing aircraft to prevent aircraft collisions. Centerline lights provide positive guidance along predetermined routes
for the pilot to follow into gates during low-visibility conditions. Stop bars consisting of rows of red lights prevent planes and vehicles from traveling onto an active runway without positive control from the air-traffic control tower. The controller has positive identification of the aircraft location through the use of ground pressure sensors and airport surface detection radar equipment. Design of these systems must mirror the airport’s operating procedures. As an example, the lighting circuit design reflects certain groupings of taxiway segments and runways used for various operating configurations (Fig. 4). A system of vehicle service roads must be provided around the perimeter of the airfield both for access to the runways and for security patrols of the perimeter fencing. Airfield security fencing with a series of access gates is monitored with patrols and, increasingly, a remote camera surveillance system. Terminal configuration. The terminal building houses the ticketing, baggage claim, and transfer
Airport engineering centerline
37.5 ft (11 m) 1.5%
37.5 ft (11 m)
35 ft. (11 m) 25 ft. (7.5 m)
3%
1 ft (0.3 m)
3-in. (7.5-cm) asphalt surface course 7-in. (18-cm) asphalt base course 5-in. (13-cm) asphalt drainage layer
17-in. (43-cm) portland cement concrete pavement 8-in. (20-cm) cement-treated base course 12-in. (31-cm) stabilized subgrade
4–6-in. (10–15-cm) drainline Fig. 3. Cross section of a runway-taxiway pavement.
low-visibility taxi route
centerline lights
Fig. 4. Surface movement and guidance and control system plan for a runway and the associated low-visibility taxi route.
to ground transportation facilities. The concourse is generally the combination of facilities for passengerboarding of aircraft, sorting baggage according to flight, and uploading cargo carried in commercial aircraft. The most popular terminals with passengers are those that offer amenities such as a hotel integrated into the terminal main area, extensive retail shops, and vibrant integrated artwork. Airport terminal and concourse configurations generally fall into three categories (Fig. 5): (1) a terminal contiguous with concourse satellite extensions (known as piers or fingers) used for boarding aircraft, such as at Miami International Airport; (2) unit terminals, which serve as transfer points both from the ground transportation modes into the building and from the building into the aircraft, such as at New York’s John F. Kennedy Airport; (3) and a detached terminal and concourses, sometimes referred to as a landside and airside terminals, connected by a people-mover train system such as at Atlanta’s Hartsfield International Airport, or an underground walkway such as at Chicago’s O’Hare Airport, or a surface transport vehicle such as at Washington’s Dulles Airport. Contiguous. The contiguous terminal-concourse has a relatively low cost since no trains or special airport vehicles are used. Passengers walk through the complex, aided by moving sidewalks. Expansion of additional aircraft gates is convenient by lengthening one of the piers or adding a pier. However, expansion is fairly limited since continued lengthening of the concourse adds walking distance and can infringe on required airfield operational space. Another limitation is the maneuverability of aircraft between the piers. Aircraft can be delayed while waiting for an aircraft to push back from an adjacent or opposite gate. Unit. Unit terminals are the most convenient for local passengers, allowing short distances from the
397
398
Airport engineering cross taxiway bridges
attached concourses
underground people mover parking and railroad station
centralized terminal
unit terminal
parking terminal
remote concourses
roadway system (a)
(b)
(c)
Fig. 5. Alternative terminal configurations. (a) Contiguous terminal and concourses. (b) Unit terminal. (c) Detached terminal and concourses.
curbside roadway system to the aircraft boarding area. This arrangement requires a more extensive investment in the curbside roadway infrastructure. Hubbing passengers, those transferring from one plane to another at the same airport, are disadvantaged in this configuration because they have to move from one building to another. Detached. Detached terminal and concourse configurations allow airfield facilities to be widely spaced for more efficient aircraft movement, while keeping passenger walk distances short by providing automatic transport systems. Landside and airside building can be expanded independently as different functions require. This configuration allows for a single security screen pavilion at the entrance to the aircraft gate area, which is an advantage for both security and cost. Shapes. A number of shapes can be used for the airside concourse, depending on the available space. The most common is a long, linear concourse with an underground train station in the middle. This concept was originally developed for Atlanta’s Hartsfield International Airport, and it proved to be the most efficient shape for aircraft movement. This linear shape mirrors the parallel pattern of aircraft movement required for highly efficient operations. The hydraulic analogy is laminar flow, as opposed to turbulent flow, which is inefficient. The master plan for Denver International Airport recognized the efficiencies of the Atlanta-style concept and adapted to it. The Denver airport’s linear concourses are more widely spaced than the Atlanta concourses, allowing simultaneous taxiing between concourses. Denver’s linear concourses were widened to allow wide moving sidewalks. An expanded central core was also added to allow for additional food and retail services for passengers. This master plan concept has now been adopted for the new terminal at Salt Lake International Airport and the new Terminal 5 at London’s Heathrow International Airport. At the Pittsburgh airport, a remote concourse has been constructed in an X shape. The large middle houses the station for the underground people mover. A great advantage is that the average passenger walk distance is halved in comparison to the lin-
ear concourse. If each concourse contains 40 gates, a passenger on the linear must walk the length of 20 gates to reach the farthest gate. The farthest a passenger in an X concourse would walk is the length of 10 gates. Gates vary in size according to the type of aircraft scheduled for use. Wider wingspans require wider (longer) gates. Aircraft with higher seat counts must have gates with a commensurate amount of passenger seating. Smaller facilities. Smaller terminal facilities are required for general aviation, charter companies, and commuter or regional aircraft. These are generally less system-intensive. Selection of locations must take into account both the number of passengers that connect from these terminals to the air carrier concourses and the preferred mode of transport between the terminals. Dimensional requirements. Airports are among the most specifically designed and configured public facilities. An important part of airport engineering is optimizing the trade-off between providing sufficient dimensional space requirements for such functions as passenger queues, aircraft movement, seating areas, and circulation without requiring overly large spaces that drive up the cost of construction and operation. A detailed dimensional requirements study is prepared that determines the width of ticket lobby, size of train stations, mix of various-sized aircraft gates, width of concourse, number of roadway lanes, number of short-term parking spaces, number and size of food service outlets, linear feet of ticket counter, and space programs for support functions; and that otherwise defines as many elements the airport as possible. Reasonable delay criteria are defined for almost every step of the trip through the airport. These criteria are then translated into speed and size requirements. Terminal subsystems. Because of the high volume of passengers and baggage in an airport terminal, mechanized subsystems are used extensively. Moving sidewalks shorten walk distances and passenger connect times. Double or extrawide moving sidewalks are desirable, provided the building is sufficiently wide, to allow some passengers to rest while others move briskly. Increasing passenger
Airport engineering preference to carry small items of luggage and shopping parcels has also required engineers to widen moving sidewalks. Two devices are used to move passengers from one level in an airport building to another. Banks of elevators are required for every level change in order to accommodate all possible passenger needs. Escalators are required in order to move the peak numbers of passengers down from ticketing to a people mover station or from a concourse down to baggage claim level. These two types of devices are complementary and are always used in tandem. Their numbers and dimensions are determined by accessing the passenger demand and required hourly capacity. One level change, for example, might require one elevator, two down escalators, and one up escalator. Use of electric trains to transport passengers within an airport building or between buildings is the most expensive alternative. Extensive structural requirements and fire code provisions add to the base system cost. Extensive signage and voice messages are required to ensure that passengers feel comfortable using such a system to move from gate to gate or concourse to concourse. This system is, however, the most comfortable for passengers. Mechanized sortation systems are used extensively at airports for mail, cargo, parts, and baggage. They significantly decrease the time to sort and transport the aircraft load. They generally use tunnels and basements to move these materials into the aircraft. Each system has special requirements (such as maximum and minimum dimensions, heavy structural loading, and maintenance access requirements) that must be considered in the building design. Security requirements. After the events of September 11, 2001, new security regulations enforced by the Transportation Security Administration have resulted in many changes to airport facilities and operational procedures. They include 100% screening of all checked baggage, matching all checked baggage to passengers, closer inspection of carry-on luggage, and prohibition of certain objects in the aircraft cabin, as well as secure cockpit doors, more rigorous background checks of employees, and more explosives detection. New airport security initiatives focus on gathering and analyzing information about passengers, cargo, vehicles, employees, and even the planes, both on the ground and in the air. An entirely new security-focused information management infrastructure is being created. Information management/information technology. In recent years, computers and computerized control systems have been developed for virtually every component of airport design and operations. Therefore, it is critical to the success of every airport project that these systems be defined and planned at the earliest stage. Airports have an unusually high number of tenants and users. Each requires a different set of systems for up-to-date information and business controls. Generally, a fiber-optic cabling system is provided for transmission of data and voice signals by all tenants. This has become an essential utility. For radio
transmission, some airports are managing the high usage of radio with another utility system, a leaky coaxial cable, which picks up low-intensity transmissions by individual radios and transmits them for amplification for the end-receiver of the message. This avoids high-intensity transmissions, which cause interference between users. Airports generally have a set of life-safety systems that include the fire alarm system, public address system, smoke evacuation, security system, and some aspects of transport systems such as underground train systems. Each of these systems must perform its role in communicating to the passengers and airport employees and controlling devices such as exhaust fans, doors, and trains. Airport operations require a paging system, card access security, controlled-circuit TV security, ground-to-ground radio, ground radar, parking and ground transportation revenue and control system, broadcast TV, and Doppler radar weather system. Some level of integration of these systems is essential to minimize the time required to communicate critical information. For instance, if a weather detection system senses a tornado or wind shear condition, the integrated control system can preprogram audible alarms. Airports also have a set of maintenance systems that are best managed by an integrated control system. These systems manage devices such as airfield in-pavement lights, pavement temperature sensors, moving sidewalk controls and alarms, irrigation controls, fuel system monitoring data, environmental monitoring data, work orders, and spare parts inventory. For engineering purposes, airports may use a computer-assisted drafting (CAD) database of all airports facilities. Information is generally stored in layers, with all structural and demising walls in one layer, utilities and electrical wiring in another, door hardware in another, and lease information in yet another. Airlines are increasingly providing their own radar to gain up-to-date information on incoming aircraft locations for better management of gates. They use both air-to-air and air-to-ground radio systems. Each airline has a master computer system that controls reservations, ticketing, and flight management. This master computer system may control other computerized systems, such as flight information displays or baggage systems. Some changes are occurring in the areas of responsibility for these various systems. Some airports are installing Common User Terminal Equipment, which provides complete passenger data-handling infrastructure to different airlines. Airports are also increasingly installing their own Flight Information Display Systems, so that multiple carriers’ flights can be displayed on the same devices. Another change is that systems historically installed and operated by the FAA may transfer to airports. One example is approach lightning systems, which indicate the distance to the runway threshold to pilots on approach. As these systems are turned over to airports,
399
400
Airport engineering another integration opportunity is presented. Integration controls for both approach lights and inpavement lights would decrease the opportunity for delay or error and allow the controller to focus on communication with pilots. A gradual change from radio technology to use of digital transmissions will affect many airport systems. Receipt of a radio transmission to a pilot cannot be verified by the air-traffic controller. However, receipt of a digitized message can be verified and can be printed in the cockpit. As older aircraft are replaced, new aircraft are equipped to receive such messages. Industry groups encourage cooperation between airlines and airport authorities. They are also encouraging development of system standards and interfaces between airline and airport authority computer systems. Recent technology developments include machine-readable boarding passes and tags, which allow better security and faster passenger processing; self-service passenger check-in kiosks; and software that allows printing of boarding passes at home. These innovations decrease passenger waiting times on lines and reduce airline personnel costs. Movement of cargo requires communication between different agencies such as airlines, customs agents, cargo shippers, cargo receivers, and trucking companies. Guidelines for standardizing these types of system are also being recommended by the industry groups. Ground transportation facilities. Airline passengers use terminal and parking facilities to transfer to private cars, commercial vans and buses, taxis, rental cars, or rail transport. Large numbers of these vehicles must be able to park, process business transactions, and load passengers. Demand for short walk distances requires that many of these functions be incorporated into the front area of the terminal building. Remote staging facilities are often constructed so that large commercial vehicles do not dwell in front of the terminal. Commercial vehicles are mobilized to the curbside by radio dispatch, loaded with passengers, and moved away from the terminal as quickly as possible, thereby maximizing the number of vehicles that can load in a peak hour, and minimizing the exposure of waiting passengers to vehicle emissions. Rail connections are best provided by having a station incorporated directly into the terminal or conveniently accessible from the terminal. Stations providing airline passengers access to regional transit are becoming more common. Airports with such facilities include New York JFK, Atlanta Hartsfield, Washington Reagan National, St. Louis Lambert, Portland (Oregon), and Baltimore-Washington. Several airports are now planning such links, including Washington Dulles, Los Angeles, and Phoenix Sky Harbor. Support buildings. The primary types of support buildings required by the airlines for their airport operations are flight kitchens to prepare meals for pas-
senger, hangars to service aircraft, and ground support equipment buildings to service vehicles such as tugs, baggage carts, and service trucks. The high number of trips for support vehicles to travel from these buildings to load or service aircraft requires that the buildings be located in reasonable proximity to the aircraft gates. However, the buildings should be sufficiently distant to allow the concourses to be expanded without requiring demolition of these support facilities. Careful forecasting of growth, analysis of operational costs, and dimensional analysis are required to obtain the optimum location. Airline companies may also require parts storage buildings, employee facilities, and administrative offices. Airlines frequently have an exclusive cargo building to support delivery and sortation of cargo carried in excess space on commercial flights. An airport requires fire equipment to provide extremely fast primary and secondary responses to every runway. Locating the aircraft rescue and firefighting stations to meet these responses times stipulated by the FAA requires careful positioning with respect to the taxiway systems. The exact building layout must recognize the need for the fire equipment doors to exit onto the most direct response route. See FIRE TECHNOLOGY. Other types of support buildings include vehicle maintenance buildings for snow removal and airport vehicles, sand storage buildings, roadway revenue plaza offices, and training facilities. Fuel and deicing facilities. Economics of scale and safety considerations generally encourage the implementation of large, centralized common systems for aircraft fuel. The large storage tanks required to ensure adequate reserves of fuel are located in remote areas of the airport, generally in aboveground facilities. Underground distribution piping transports the fuel to hydrant pits or truck fueling stations close to aircraft operations. This system, like most utilities, is designed with backup capacity by looping piping around each service area. If a problem occurs in one section of pipe, valves are automatically closed and the supply direction is reversed. Fuel tanks require extensive structural, mechanical, and electrical design. These tanks are widely spaced to avoid the transmission of fire and to allow room for a surface detention area to store burning fuel. Airlines use either propylene glycol or ethylene glycol, mixed with water, to deice aircraft and to provide coatings that retard development of ice accumulation on aircraft. Facilities are required to store the glycol, mix it, and sometimes heat it before distribution. It can then be loaded onto deicing trucks for spraying onto the aircraft or distributed to deicing gantries or booms that are fixed to a base but have an extendable arm with an operator cab on top from which to spray the fluid. Glycols have serious water-quality impacts on receiving streams and treatment plants. Increasingly, airports are making provisions to grade the ramp areas where deicing operations occur so that the spent glycol can be captured. Once it is segregated from the normal storm collection system, it can be
Airport noise stored for metering in treatments facilities, reprocessed for reuse off airport, or otherwise disposed of. Ginger Sunday Evans Bibliography. A. Graham, Managing Airports: An International Perspective, 2d ed., 2003; R. M. Horonjeff, F. X. McKelvey, and B. Sproule, Planning and Design of Airports, 5th ed., 2007; R. de Neufville and A. Odoni, Airport Systems: Planning, Design, and Management, 2002; A. T. Wells, Airport Planning and Management, 4th ed., 2000.
Airport noise The unwanted sound from airport operations, primarily from aircraft of all types. It affects neighbors both adjacent to and farther from the airport. It can be controlled using several approaches, and assessed using a number of tools. Airport traffic and construction, which also create noise, will not be discussed here. Types of operations. Aircraft using an airport cause noise during several types of operations. The noisiest is when an aircraft takes off. The plane is fully loaded with passengers, cargo, and fuel and uses maximum thrust to take off from the runway. Neighbors behind and to the side of the runway experience the first noise from a take-off. At later times, people along the take-off path experience the noise from the plane as it passes overhead. When the plane approaches its destination airport, it again causes noise for the people along its path. It is now much lighter, having used up fuel enroute. However, at United States airports, the planes are also much lower as they line up to land on the designated runway, coming in on a 3◦ glide slope. When the aircraft lands, the pilot applies the brakes and/or reverses the engines. This operation also causes noise for the neighbors adjacent to the airport. Aircraft also produce noise when they taxi to or from the runways, gates, parking areas, or maintenance areas. In this case, the amount of thrust needed is much smaller than the amount needed to be airborne, and the associated noise is much less. Finally, aircraft can cause noise when they perform maintenance run-ups. After aircraft maintenance, the Federal Aviation Administration (FAA) requires that the system involved be tested before flight. This may require running the engines to partial or full thrust. The run-up is done on airport property but without the plane actually taking off. A full-power run-up is nearly as noisy as a take-off. Mitigation. Noise control can be applied at the source, along the sound path, or at the receiver. In the case of airport noise, the source is the aircraft and the receiver is the neighbor. Noise control at the source is preferred, since it reduces the total amount of sound released into the environment. Source noise control of jets has progressed with time. The first jet aircraft, called stage 1, had no noise control on their engines. Federal law now prohibits use of stage 1 in the United States. Simple muffling was added to the jet engines in the 1970s
to give stage 2 aircraft. At the same time, improvements were made to the engine design, resulting in about 10 decibels of noise reduction for a given aircraft type. This type of engine powers the stage 3 aircraft. As of January 1, 2000, all large turbojet aircraft operating in the United States are stage 3, by federal mandate. In anticipation of the stage 3 mandate, owners of stage 2 aircraft could sometimes use an engine retrofit, called a hushkit, to meet the stage 3 noise criteria, prolonging the useful life of the aircraft. Additional improvements, starting in 1992, have continued to reduce aircraft noise emissions. In June 2001, the International Civil Organization (ICAO) adopted the next quieter level, called Chapter 4 of their regulations (Annex 16). It effectively eliminates hushkitted older aircraft from operating in countries which adopt these rules. See AIRCRAFT ENGINE PERFORMANCE; JET PROPULSION; MUFFLER; TURBOJET. Airports can reduce the noise from backing up, repositioning, or taxiing by having the aircraft towed to or from the gate. The noise from towing is much less than that from the aircraft engines. In addition, aircraft may be asked to use only the engine away from the neighbors, to cut the noise source level by 3–6 dB depending on the number of engines on the aircraft. Finally, aircraft can be connected to airport electrical service at the gates instead of running their auxiliary power units or they can use ground power units, especially at night. The source noise during flight can be reduced by using lower thrust settings when planes fly over noise-sensitive areas. So, if neighbors are far from the airport, the aircraft can use maximum thrust to get as high as possible and lower thrust farther away. This method, encouraged by the FAA, is used at a number of United States airports. The most common path noise control is by use of barriers or berms. For effective noise reduction (at least 5 dB), the top of the barrier must intercept the line-of-sight between the aircraft engine and the observer. Since the elevation of aircraft engines can be as high as 30 ft (9.1 m), noise reduction occurs only when the neighbors are close to the ground, when the aircraft are on the ground, and when the barriers are high. Noise from taxiing and run-ups can be reduced in this way. At some airports, buildings are constructed near the active taxiways to provide effective barriers. A second type of path noise control is to increase the distance between the source and receiver. The sound level is decreased 6 dB for every doubling of distance. Airports, working with the FAA, can develop noise abatement departure or arrival paths to maximize the distance between the source and receiver or to take advantage of open space. The airport can also move or “buy out” the affected residents. The area left vacant can either be redeveloped as compatible land use, which may be commercial or industrial, or left undeveloped. Path noise control of aircraft run-ups is achieved by use of a hush-house, a hangarlike structure in which the aircraft can perform its engine test. This partially
401
402
Airport surface detection equipment or fully enclosed structure reduces the amount of noise emitted to the environment. These facilities are used at military facilities, overseas, and at a small number of United States commercial airports. Receiver noise control is usually by soundinsulating the structure. The windows are replaced with double or triple glazing, and doors may be replaced with solid core versions. Typical improvement in noise reduction is 5–25 dB depending on the original condition of the structure. Air conditioning may also be added. This approach, which is effective only when the neighbor is inside the structure, has been used in the United States since 1980. Time-of-day considerations. Most airports receive complaints from their neighbors about aircraft noise. The complaints include not being able to hear the television, constant noise, sleep deprivation, inability to work, and vibrations. United States airports must (by federal law) be open 24 hours per day, 7 days per week, unless there is another facility which can take the traffic. Complainants may also be affected by through flights or other airports’ operations. To mitigate this situation, an airport can direct nighttime operations to occur on a runway which affects the minimum number of people, whenever operationally possible. Limitations can also be placed on runways which affect many neighbors or are very close to them, if other runways are available. Altitude restrictions may be added at night to maximize the distance between the aircraft and the neighbors. At major airports, where the maximum hourly number of arrivals and departures may be larger than 100, there will be several runways which can be used for aircraft operations. In this case, an airport can request that the FAA alternate among them, consistent with operational requirements, to relieve communities from constant noise. Assessment. Airport noise is measured with noise monitors, consisting of microphones and meters. Most large airports have a distributed network of permanent noise monitors which record observed noise levels and accumulate the data for future transfer to the airport’s computer. Small airports are most likely to have only portable measuring equipment or to use consultants. See NOISE MEASUREMENT. Observations of aircraft operations are either direct or automated. With an automated system, an airport may have access to the FAA’s data or to passive radar data. These data give the exact location, identity, and time of every operation in the area. Direct observations of aircraft operations are very reliable but are labor-intensive. The radar data can be used to determine which operation affected a given resident at a given date and time. It can also be accumulated to determine the number of operations in a year by aircraft type on each of the airport’s runways by time of day. The information can be used by the airport for assessing conformance with established noise abatement paths and procedures and for verifying operational reports from the airlines. Combined with information on airline engine
types and typical performance data on the aircraft type, the radar data can be used in airport noise prediction models, such as the FAA’s Integrated Noise Model, to determine contours of noise exposure. These contours are averages of day-night equivalent sound levels for the year and can be compared with measured averages of the same type for locations for which the airport has noise monitors. These contours, in turn, can be used by the airport for planning purposes or for obtaining federal funds for noise mitigation programs. If an airport has regulations, the above information may be used for enforcement. If the regulations are noise-level-based, microphone data will be used for this purpose. At airports where the regulations are operations-based, observations by airport personnel are often used to ensure compliance with local regulations. Complaints from the community may assist an airport in identifying potential violations. The FAA also assists the airport in ensuring compliance with noise regulations. In deciding to adopt new noise mitigation at an airport, the FAA’s regulations in Part 150 of the Code of Federal Regulations govern. A study of existing noise impacts is prepared, and mitigation strategies such as those discussed above are proposed. The FAA reviews the document and may approve it for federal funding of the proposed projects. Any new regulations would also need to be approved by the FAA. In addition, many states have environmental reporting requirements for their airports. Measurements of the noise impacts and new initiatives are properly filed with the appropriate authority. See ACOUSTIC NOISE. Nancy S. Timmerman Bibliography. M. J. Crocker (ed.), Encyclopedia of Acoustics, Wiley, New York, 1997; C. M. Harris, Handbook of Acoustical Measurements, 3d ed., Acoustical Society of America, Sewickley, PA, 1998; C. M. Harris, Handbook of Noise Control, 2d ed., McGraw-Hill, New York, 1979; H. H. Hubbard, Aeroacoustics of Flight Vehicles: Theory and Practice, Acoustical Society of America, Woodbury, NY, 1995.
Airport surface detection equipment A ground mapping system that uses analog radar equipment to provide surveillance of aircraft and other surface vehicles on an airport surface. It is used by air-traffic controllers in airport control towers to monitor and control the movement of aircraft and vehicles. A situation display of the targets includes a map identifying the runways and taxiways and a visual map of the airport features, created through the contrast on the radar display resulting from the absence of returns from smooth concrete surfaces and the ground-clutter returns from grassy areas. An important safety function of the airport surface detection equipment (ASDE) is to determine whether or not a runway is clear for the next departure or arrival operation. The prevention of runway incursions by aircraft or by service vehicles has become
Airship more urgent because of several incidents since 2003. This runway clearance determination is aided by the ASDE’s capability to display an image of the aircraft in which the target’s extremities (nose, tails, and wing tips), especially for large aircraft, are evident to the eye. See AIRPORT. Operation. The ASDE antenna revolves at 60 revolutions per minute, providing a rapid update of target movements. It is located on a high vantage point, typically on top of the tower cab or on a special remote tower, that provides line-of-sight coverage of the desired runway and taxiway areas. Older ASDE systems have several inherent weaknesses, especially equipment failures, and excessive clutter during precipitation conditions, and are typically operated at a fixed frequency (for example, the ASDE-2 operates at 24 GHz). Modern ASDEs utilize digital processing and an interrogation technique called frequency agility, which permits the transmission of up to 16 different Kuband frequencies within a cycle. These techniques optimize target detection in the presence of ground and rain clutter and the rejection of false targets. In the digital processing, different thresholds for rejecting clutter are used on specified areas of the airport surface, such as on paved surfaces versus grassy areas. These systems are designed to detect all aircraft and vehicles having a radar cross section of 30 ft2 (3 m2) or greater, and will resolve two closely spaced targets when separated by 40 ft (12 m) in range or 80 ft (24 m) in azimuth at a range of 12,000 ft (3650 m). The basic coverage is the entire airport surface out to 24,000 ft (7300 m) in range and 200 ft (60 m) in altitude. Areas where coverage is not de-
sired are ignored for detection and blanked out on the situation display. Automation functions. The digital processing in modern ASDEs also provides the basis for automation functions that process target data to provide controllers with improved information. Principal automation features include target tracking, the application of runway safety logic functions to the tracks, and subsequent automatic alerting of the controller to dangerous situations that are determined by the logic. Advanced airport surface surveillance systems will enhance the data already provided by the ASDE, permitting more advanced automation functions. The design requirements for these systems include providing data that allows determination of target identity and target intent from flight-plan information or from data-link messages between the control tower and aircraft. See ELECTRONIC NAVIGATION SYSTEMS; MOVING-TARGET INDICATION; RADAR. Robert A. Bales
Airship A propelled and steered aerial vehicle, dependent on the displacement of air by a lighter gas for lift. An airship, or dirigible balloon, is composed primarily of a streamlined hull, usually a prolate ellipsoid which contains one or more gas cells, fixed and movable tail surfaces for control, a car or cabin for the crew or passengers, and a propulsion system. Past designs. Two fundamentally different designs have been successfully used for past airships, the nonrigid and the rigid. A third type, the semirigid, is
Hindenburg Graf Zeppelin Bodensee
Viktoria - Luise
R-100
(a)
Gr ea tB
y an rm e G
Schwarz #2
ain rit
nce Fra
RS - 1
M
D ited Un
Giffard
Norge
tes Sta ed t i Un
Santos-Dumont nce Fra
ZMC-2 ZR-1 Shenandoah
Italy
Lebaudy
Akron and Macon
s tate
LZ -1
(b)
(c)
S ited Un
es Stat
ZPN
ZPG - 3W
Pony blimp
Baldwin 1850
1900
1910
1920
1930
1940
year Fig. 1. History of airship development. (a) Rigid airships. (b) Semirigid airships. (c) Nonrigid airships.
1950
1960
403
404
Airship
catenary curtain nose cone supports aft ballonet air helium forward ballonet air
suspension cables
car
engine
Fig. 2. Typical nonrigid airship.
fin longitudinal girders intermediate ring (frames) gas cell
fin
engine car
main ring (frames)
bow mooring cone control car Fig. 3. Typical rigid airship.
Fig. 4. Akron, rigid airship of the U.S. Navy.
essentially a variant of the nonrigid type, differing by the addition of rigid keel. The principal development trends of the three types of conventional airships are depicted in Fig. 1. Nonrigid. A typical nonrigid airship, or blimp (Fig. 2), consists of a flexible envelope, usually fabric, filled with lifting gas that is slightly pressurized. Internal air compartments (ballonets) expand and contract to maintain the pressure in the envelope as atmospheric pressure and temperature vary. Ballonet volume is controlled by ducting air from the prop wash or by electric blowers. The weights of the car structure, propulsion system, and other concentrated loads are supported by catenary systems attached to the envelope. The nonrigid airships are historically significant for two reasons. First, a nonrigid airship was the first aircraft of any type to achieve controllable flight, in the 1850s. Second, nonrigid airships were the last type to be used on an extensive operational basis; the U.S. Navy decommissioned the last of its nonrigid airship fleet in the early 1960s. During the many years the Navy operated them, a high degree of availability and reliability was achieved. Most of these nonrigid airships were built by Goodyear, and a few, based on a modified pre-World War II Goodyear design, are used today for advertising by that company. See BLIMP. Rigid. The other major type of airship is classified rigid because of its structure (Fig. 3). This structure was usually an aluminum ring-and-girder frame. An outer covering was attached to the frame to provide a suitable aerodynamic surface. Several gas cells were arrayed longitudinally within the frame. These cells were free to expand and contract, thereby allowing for pressure and temperature variations. Thus, despite their nearly identical outward appearance, rigid and nonrigid airships were significantly different in their construction and operation. The rigid airship was developed primarily by the Zeppelin Company of Germany, and, in fact, rigid airships became known as zeppelins. Even the small percentage of rigid airships not built by this company was based, for the most part, on the Zeppelin concept. The rigid airships of the Zeppelin Company recorded some historical firsts in air transportation, including inaugurating the first passenger air service. The culmination of zeppelin development were the Graf Zeppelin and Hindenburg airships—unquestionably outstanding engineering achievements in their day. All of the rigid airships produced in the United States were for military purposes; none was in operation at the outbreak of World War II. The last rigid airships built and operated in the United States were the Akron (Fig. 4) and Macon, built in 1931 and 1933. These airships were in many respects typical of the large rigid type but also had some unique design features. They contained just under 7,000,000 ft3 (200,000 m3) of helium, were a little over 800 ft (240 m) long, and had a gross lifting capability of about 450,000 lb (200,000 kg). Among the unique features of these airships were their eight engines. The propellers could be swiveled
Airspeed indicator in any vertical or fore and aft direction, which made the Akron and Macon two of the first aircraft to use vectored thrust propulsion. The panels on the side of the airship that looked like windows were actually condensers which collected water vapor out of the engine exhaust to replace weight on board the airship as fuel was burned. The internal rigid structure of the Akron and Macon was made of aluminum alloy. In fact, the use of aluminum alloy in aerospace structures was pioneered by rigid airships. The role of the Akron and Macon was to serve as flying aircraft carriers and scout vehicles. They carried five small airplanes on board to extend their range of surveillance, protect the airship, and serve as scouts for the surface fleet. Commercial operations. The only significant past commercial airship operations were those of the Zeppelin Company and its subsidiary DELAG. This organization began carrying passengers on flights within Germany in 1910, and from 1929 to 1937 provided a unique transatlantic air service with the Graf Zeppelin and Hindenburg. Prior to the burning of the Hindenburg in 1937, there had been no fatal accidents in commercial service. However, none of these commercial operations can be considered a financial success, and most were subsidized by the German government. Modern vehicle concepts. Figure 5 shows a representative sample of the many modern airship vehicle concepts which have been proposed. The fully buoyant conventional concepts are modern versions of the classical, fully buoyant, rigid and nonrigid airship concepts. Fully buoyant means all the lift is provided by displacement. These airships would make extensive use of modern aircraft structural materials, propulsion systems, control systems, and electronics. The other four concepts in Fig. 5 are partially buoyant, or hybrid, designs in which the buoyant lift is substantially augmented by aerodynamic or propulsive lift. Thus, these vehicles are partly heavier than air and partly lighter than air. Two of the hybrids are short takeoff and landing (STOL) vehicles, and two have vertical takeoff and landing (VTOL) capability. Generally speaking, the partially buoyant concepts would have higher cruise speed, higher fuel consumption, and higher structural weight as compared with the classical concepts. It will therefore depend on the specific applications as to which is the best concept. An important advantage of the hybrids is that they promise to alleviate the costly ground-handling requirement of past airship designs. See SHORT TAKEOFF AND LANDING (STOL); VERTICAL TAKEOFF AND LANDING (VTOL). The Dynairship is representative of lifting-body concepts in which the hull is shaped to be more aerodynamically efficient than the conventional airship shape. This particular concept has a delta platform. The vehicle is flown at positive angle of attack in cruise to generate aerodynamic lift. The Megalifter concept is a typical example of a winged airship. The Helipsoid is intended to be a compromise between the ellipsoidal and deltoid shapes. It will have a better
rigid airship
Aereon Dynairship lifting body concept
megalifter (with turboprops) fixed-wing concept
nonrigid airship
Fig. 5. Modern airship concepts.
surface-area/volume ratio than the Dynairship, at the expense of the degraded stability and control characteristics. Finally, the Heli-stat, which is representative of the rotary-wing hybrids, combines an airship hull with helicopters or helicopter rotor systems. See BALLOON. Mark R. Ardema
Airspeed indicator A device that computes and displays speed of an aircraft relative to the air mass in which the aircraft is flying. Three common types are the indicated airspeed meter, true airspeed indicator, and Machmeter. The commonest type is the indicated airspeed meter, which measures differential pressure between the total ram pressure from the Pitot system and the total static pressure; it then converts this difference into units of speed (miles per hour or knots) under standard conditions, 29.9212 in. Hg (101.325 kilopascals) absolute and 59◦F (15◦C; see illus.). Although the indicated values are incorrect above zero altitude, the relationship to the aircraft handling remains essentially unchanged, thus providing a measure of the flyability of the aircraft. These instruments
Pitot connection
airspeed diaphragm
indicated airspeed pointer
rocking shaft
maximum allowable pointer
aneroid diaphragm
Components of indicated airspeed meter.
405
Boeing helipsoid lifting body concept
Piasecki Heli-stat rotary-wing concept
406
Aistopoda are frequently supplemented with an auxiliary mechanism, driving a second pointer which computes and indicates a value known as the maximum allowable speed for the type of aircraft in which the instrument is installed. True airspeed indicators are similar to indicated airspeed meters but include a more complex mechanism that also senses both the absolute pressure and temperature, and compensates for the change of density of the air mass, thus obtaining true airspeed. This indication is of value in computing course information; hence the true airspeed indicator is a navigation instrument whereas the indicated airspeed indicator is a flight instrument. For those aircraft that reach higher speeds (transonic and supersonic), the ratio of the actual speed to the local speed of sound is used. Devices that compute this value are known as Machmeters. Relative Mach number is computed by determining the ratio of total ram pressure divided by total static pressure and converting the ratio into Mach number, according to the applicable formulas. Certain instruments combine airspeed, Mach number, and maximum allowable speed in one indicator to provide a central speed-data display. See MACH NUMBER; PITOT TUBE. James W. Angus
Albedo A term referring to the reflecting properties of a surface. White surfaces have albedos close to 1; black surfaces have albedos close to 0. Bond albedo and normal reflectance. Several types of albedos are in common use. The Bond albedo (AB) determines the energy balance of a planet or satellite and is defined as the fraction of the total incident solar energy that the planet or satellite reflects to space. The “normal albedo” of a surface, more properly called the normal reflectance (rn), is a measure of the relative brightness of the surface when viewed and illuminated vertically. Such measurements are referred to a perfectly white Lambert surface—a surface which absorbs no light and scatters the incident energy isotropically—usually approximated by magnesium oxide (MgO), barium sulfate (BaSO4), or some other bright material. See PHOTOMETRY; PLANET. Bond albedos for solar system objects range from 0.9 for Saturn’s icy satellite Enceladus and Neptune’s Triton to values as low as 0.01–0.02 for dark objects such as the satellites of Mars (Table 1). Cloudshrouded Venus has the highest Bond albedo of any planet (0.76). The value for Earth is typically 0.35, but varies with the extent of cloud and snow cover. The Bond albedo is defined over all wavelengths, and its
Aistopoda An order of extremely elongate, limbless fossil amphibians in the subclass Lepospondyli from PermoCarboniferous rocks of North America and the British Isles. The order includes three families: Lethiscidae (Lethiscus), Ophiderpetontidae (Coloraderpeton, Ophiderpeton), and Phlegethontiidae (Aornerpeton, Phlegethontia, Sillerpeton). All genera are monotypic with the exception of Ophiderpeton (six species) and Phlegethontia (possibly three species). In seeming contrast to their derived state, aistopods are among the very oldest of all known fossil tetrapods. The skulls of aistopods are fenestrated and exhibit a reduced number of bony elements. This fenestration is extreme in phlegethontiids, whose skulls may have been kinetic. The vertebral centra are holospondylous (singlepieced), hourglass-shaped, and fused to their neural arches. Vertebrae can exceed 200 in number and frequently bear foramina for passage of spinal nerves. Ribs typically bear a unique process, which gives them a K-shape in some species. Lethiscids and ophiderpetontids possess a ventral armor of fusiform, bony elements arranged in a chevron pattern. Such armor in phlegethontiids consists of more slender, threadlike elements. Most aistopods were presumably aquatic, although rib specializations, and the more gracile proportions of phlegethontiids, suggest a rather snakelike, terrestrial habit. See AMPHIBIA; LEPOSPONDYLI. Carl F. Wellstead Bibliography. R. L. Carroll, Vertebrate Paleontology and Evolution, 1988.
TABLE 1. Bond albedos and visual geometric albedos of selected solar system objects Bond albedo (AB)
Visual geometric albedo ( pv)
Mercury Venus Earth Mars Jupiter Saturn Uranus Neptune Pluto
0.12 0.76 0.35 0.24 0.34 0.34 0.34 0.28 0.5(?)
0.14 0.59 0.37 0.15 0.45 0.46 0.48 0.50 0.61
Moon Phobos (M1) Deimos (M2) lo (J1) Europa (J2) Ganymede (J3) Callisto (J4) Amalthea (J5) Mimas (S1) Enceladus (S2) Tethys (S3) Dione (S4) Rhea (S5) Titan (S6) Ariel (U1) Umbriel (U2) Titania (U3) Oberon (U4) Miranda (U5) Triton (N1) Nereid (N2) Proteus (N8)
0.12 0.02 0.02 0.56 0.58 0.38 0.13 0.02 0.60 0.90 0.60 0.45 0.45 0.20 0.21 0.10 0.15 0.12 0.18 0.90 0.07 0.03
0.14 0.05 0.06 0.63 0.68 0.43 0.17 0.06 0.75 1.20 0.80 0.55 0.65 0.20 0.39 0.21 0.27 0.23 0.32 0.75 0.14 0.06
Ceres Vesta
0.03 0.12
0.06 0.23
Comet Halley (nucleus)
0.01
0.03
Object
Albite TABLE 2. Normal reflectances of materials∗
∗
Material
Albedo
Lampblack Charcoal Carbonaceous meteorites Volcanic cinders Basalt Iron meteorites Chondritic meteorites Granite Olivine Quartz Pumice Snow Sulfur Magnesium oxide
0.02 0.04 0.05 0.06 0.10 0.18 0.29 0.35 0.40 0.54 0.57 0.70 0.85 1.00
Powders; for wavelengths near 0.5 micrometer.
value therefore depends on the spectrum of the incident radiation. For objects in the outer solar system not yet visited by spacecraft (such as Pluto), the values of AB in Table 1 are estimates derived indirectly, since for these bodies it is impossible to measure the scattered radiation in all directions from Earth. Normal reflectances of some common materials are listed in Table 2. The normal reflectances of many materials are strongly dependent on wavelength, a fact that is commonly used in planetary science to infer the composition of surfaces remotely. While the Bond albedo cannot exceed unity, the normal reflectance of a surface can do so if the material is more backscattering at opposition than the reference surface. Geometric albedo. In the case of solar system objects, a third type of albedo, the geometric albedo ( p), is commonly defined. It is the ratio of incident sunlight reflected in the backscattering direction (zero phase angle or opposition) by the object, to that which would be reflected by a circular disk of the same size but covered with a perfectly white Lambert surface. Objects like the Moon, covered with dark, highly textured surfaces, show uniformly bright disks at opposition (no limb darkening); for these, p = rn. At the other extreme, a planet covered with a Lambert-like visible surface (frost or bright cloud) will be limb-darkened as the cosine of the incidence angle at opposition, and p = 2/3rn. Table 1 lists the visual geometric albedo pv, which is the geometric albedo at a wavelength of 550 nanometers. Phase integral. For solar system objects, the ratio AB/p is called the phase integral and is denoted by q. Here p is the value of the geometric albedo averaged over the incident spectrum, and does not generally equal pv given in Table 1. Values of q range from near 1.5 for cloud-covered objects (1.4 for Saturn and 1.3 for Venus) to 0.2 for the very dark and rugged satellites of Mars: Phobos and Deimos. The value for the Moon is about 0.6. Other albedo types. The quantities AB, p, and q quantify the reflection of light averaged over an entire hemisphere of a planet. They are most appropriately used when details of the distribution of the reflectance on a planet’s surface are beyond the res-
olution of the observations. It is possible to define other types of albedos based on the geometry and nature (diffuse or collimated) of the incident beam and scattered beams. For such bidirectional albedos or reflectances, the angles of incidence and scattering as well as the angle between these two directions (phase angle) must be specified. Finally, in determining the energy balance of surfaces, one is often interested in the analog for a flat surface of the Bond albedo, that is, the fraction of the incident energy that the surface reflects. For a plane parallel beam incident at an angle i, this quantity, A(i), is generally close to the value AB, but can be significantly larger than AB for some surfaces, for large values of i. Role in planetary astronomy. The role that these albedos play in planetary astronomy is illustrated by the discovery of a distant new planet in the solar system. The distance of the planet from the Earth and Sun is accurately known from its motion; determination of its orbit requires only a small number of measurements of its position relative to the stars over a few months. The total brightness of the planet is the product of its area and (geometric) albedo. For a distant or small object whose size cannot be resolved with even the largest telescopes, an albedo may be assumed, allowing its size to be estimated from its brightness alone. What is known about the planet can be improved if thermal radiation from the planet is detected at infrared wavelengths. The quantity of incident solar energy absorbed by a planet is (1 − AB), that is, energy that is not reflected is absorbed. In this way, the Bond albedo determines the temperature of a planet. The absorbed sunlight heats the planet until it reaches an equilibrium temperature where the total energy it radiates at infrared wavelengths equals the total solar energy absorbed. The infrared brightness depends on the area of the planet and its temperature, so if both the infrared (thermal emission) and visual brightness (reflected sunlight) are measured, both the size and the albedo of a distant unresolved planet can be determined unambiguously. Using this technique, the sizes and albedos of thousands of asteroids and other objects that populate the solar system are now known. See HEAT RADIATION; INFRARED ASTRONOMY; RADIOMETRY. Joseph Veverka; Jay Goguen Bibliography. J. K. Beatty, C. C. Petersen, and A. Chaikin, The New Solar System, 4th ed., Sky Publishing, Cambridge, MA, 1999; B. Hapke, Theory of Reflectance and Emittance Spectroscopy, Cambridge University Press, 1993; M. Harwit, Astrophysical Concepts, 4th ed., Springer, 2006.
Albite A sodium-rich plagioclase feldspar mineral whose composition extends over the range Ab100An0 to Ab90An10, where Ab (= albite) is NaAlSi3O8 and An (= anorthite) is CaAl2Si2O8 (see illus.). Albite occurs in crustal igneous rocks as a major component of pegmatites and granites, in association with quartz, mica (usually muscovite), and potassium
407
408
Albuliformes temperatures, Al and Si become completely ordered, forming what is called low albite. If the An component exceeds 2–3 mole %, this feldspar exsolves into two distinct, submicroscopic phases, one of nearly pure albite (An0−1) composition, and the other of composition between An15 and An25; the resultant lamellar intergrowth of these two phases is called peristerite, with optical interference colors resembling those on the neck feathers of a pigeon. Such material, especially that from Hybla, Ontario, Canada, is valued as a semiprecious gemstone and may be cut and polished into cabochons. Albite is used in ceramic materials, as a source of soda (Na2O) and alumina (Al2O3) in glasses, and in certain, relatively soft abrasives. See FELDSPAR; GEM; IGNEOUS ROCKS; PEGMATITE. Paul H. Ribbe
(a)
(b) Albite. (a) Crystals, Amelia Court House, Virginia (specimen from Department of Geology, Bryn Mawr College). (b) Crystal habits (after D. S. Hurlbut, Jr., Dana’s Manual of Mineralogy, 17th ed., Wiley, 1959)
feldspar (orthoclase or microcline). Sodium and potassium feldspars usually occur as distinct mineral grains, sizes varying from millimeter to meter scale. However, they are frequently intergrown, having exsolved from a single phase at high temperatures. If the intergrowth is visually observable in a hand specimen, the composite material is known as macroperthite; if visible only in a microscope, microperthite; and if submicroscopic in scale, cryptoperthite. In metamorphic rocks albite is found in granitic gneisses, and it may be the principal component of arkose, a feldspar-dominant, sedimentary rock. Cleavelandite, a platy variety, is sometimes found in lithium-rich pegmatites. See ARKOSE; GNEISS; PERTHITE. Its hardness is 6 on the Mohs scale, specific gravity 2.62, and melting point 1118◦C (2044◦F). If observed at temperatures above 985◦C (1805◦F), albite has monoclinic symmetry, provided that the distribution of aluminum (Al) and silicon (Si) among the tetrahedral sites of its framework structure is random (disordered). Below 958◦C (1805◦F) the symmetry becomes triclinic as the [AlSi3O8]∞1− tetrahedral framework collapses around the small sodium atom, primarily through the mechanism of intertetrahedral angle bending. With annealing at lower
Albula vulpes. (Courtesy of C. Q. Jessen)
Albuliformes An order of actinopterygian fishes in the subdivision Elopomorpha, along with Elopiformes, Anguilliformes, and Saccopharyngiformes, all of which have a leptocephalous larval stage. Albuliformes, comprising three families, eight genera, and 30 species, are distinguished from other Elopomorpha by an open mandibular sensory canal in the dentary and angular bones. Two suborders are recognized, Albuloidei and Notacanthoidei, the latter of which was recognized previously as the order Notacanthiformes. See ACTINOPTERYGII; ANGUILLIFORMES; ELOPIFORMES; NOTACANTHOIDEI. Albuloidei, consisting of the family Albulidae, are characterized by a fusiform body; the tip of the snout overhanging an inferior mouth; a mouth bordered by premaxillae; a gular plate having a thin median splint or being entirely absent; an infraorbital lateral-line canal extending onto the premaxillae; and caudal fin rays supported by six hypural bones. Albulidae contains two subfamilies, Albulinae with one genus (Albula) and three species and Pterothrissinae with one genus (Istieus) and two species. Albula (see illustration) differs from Istieus (formerly Pterothrissus) in having 16–21 dorsal fin rays versus 55–65; a small gular plate versus no gular plate; and the maxillae toothless versus maxillae toothed. Albula occurs in tropical seas, rarely entering brackish and fresh water, and attains a length of 105 cm (41 in.).
Albumin Structure. A major breakthrough in albumin research came in the early 1990s, when x-ray diffraction studies of crystalline human albumin provided the first accurate images of this remarkable protein. The high-resolution structure revealed an overall heart shape, both in the defatted form and with bound fatty acids (see illustration). The protein comprises three domains, each with two subdomains, as predicted by many studies. The fatty acids are bound throughout the protein in structurally distinct sites with similar features: the carboxyl group of each fatty acid forms a salt bridge or a hydrogen bond with basic and polar amino acid side chains, and the hydrocarbon chain of the fatty acid fits into a hydrophobic cavity between helices. Physiological function. In human plasma, albumin is the major buffer for pH variations and is essential for maintaining the plasma colloidal osmotic pressure, causing fluid to remain in the blood instead of leaking out into the tissues. Albumin is the main transport vehicle for fatty acids for utilization in various tissues. As albumin circulates through the vascular system, fatty acids desorb rapidly from the protein in spite of their high affinity for the binding sites (see illustration) and diffuse through the adjacent cell membranes to enter cells. Albumin also binds fatty acids that are released into the plasma from cells such as fat cells for transport to other tissues, where the fatty acids serve as fuel or substrates for lipid synthesis. See BLOOD; LIPID METABOLISM. Role in disease. Low plasma levels of albumin are associated with several diseases, such as diabetes, liver disease, and cardiovascular disease. Undetectably low or very low plasma levels of albumin are found in an extremely rare inherited disorder known as analbuminemia. There is inadequate compensation for albumin’s role in maintaining plasma colloidal osmotic pressure, which is very low in affected individuals; however, the plasma lipoproteins
Istieus occurs in the eastern Atlantic from Mauritania to Namibia and in the western Pacific from China and Japan, and probably does not exceed 40 cm (16 in.) in length. Albula can tolerate low oxygen tension by “inhaling” air into its lunglike air bladder. It is benthic in habitat and feeds on invertebrates. Istieus is also benthic but usually occurs in deeper waters of continental slopes. Herbert Boschung Bibliography. J. S. Nelson, Fishes of the World, 4th ed., Wiley, New York, 2006; D. G. Smith, Albulidae (pp. 683–684), Halosauridae (685–687), Notacanthidae (688–689), in K. E. Carpenter (ed.), The Living Marine Resources of the Western Central Atlantic, FAO Species Identification Guide for Fishery Purposes, vol. 2, FAO, Rome, 2002.
Albumin A plasma protein produced by the liver that maintains fluid balance in the blood and transports fatty acids in the plasma and interstitial fluid. It is the most abundant protein in human serum, and one of the first discovered and earliest studied proteins. Serum albumin was precipitated from urine in 1500 by Paracelsus, and was crystallized in 1894 by A. G¨ urber. Probably no other protein has been studied as extensively as serum albumin, and our knowledge of its structure and interactions with its ligands has come from many researchers, using a great variety of experimental approaches. Its ability to bind many different ligands, most of which are hydrophobic anions, and several molecules of the same ligand (fatty acid) is well documented. Fatty acids, bilirubin (an orangeyellow bile pigment formed in the breakdown of red blood cells), and hematin (a blue to blackishbrown compound formed in the decomposition of hemoglobin) represent the endogenous ligands of albumin with highest affinity. See PROTEIN.
IB
IIIB
1 5 IA 4 7
2
IIIA IIA
3 6
IIB (a)
(b)
Human serum albumin. (a) Domain structure of the defatted protein and (b) location of eight myristic acid (a 14-carbon fatty acid) ligands in different binding sites. In a, the secondary structure of the protein is shown schematically, and the domains are designated I, II, III. The subdomains are A and B. Fatty acids with different chain lengths are present in the plasma, and the binding affinity increases with increasing chain length of the fatty acid, up to a length of 20 carbons. In b, showing the binding of the myristic acid, the atom type is differently shaded to show carbon and oxygen. Myristic acid in site 4 is darker to distinguish it from myristic acid in site 3.
409
410
Alcohol take over the role of transport of fatty acids and bilirubin. Analbuminemia is not a lethal disorder, although edema, fatigue, and lipodystrophy (a disturbance of fat metabolism in which the subcutaneous fat disappears over some regions of the body but is unaffected in others) are common complaints. The high concentration of albumin in the general population, therefore, is not essential for life but appears to be important for optimal health. Drug transport. Human serum allumin (HSA) has two principal drug-binding sites (sites I and II in the classical literature), which have been localized in the crystal in subdomains IIA and IIIA (see illustration). Site I binds bulky heterocyclic molecules with a delocalized negative charge in the center of a predominantly hydrophobic structure, including warfarin, phenylbutazone, thyroxine, benzodiazepines, antidiabetic agents, and tryptophan. Site II binds aromatic carboxylic acids that are ionized at physiological pH and contain a negative charge at one end of the molecule, including propofol, diazepam, ibuprofen, acetylsalicylic acid (aspirin), bilirubin, and thyroxine. Solubilization of these hydrophobic compounds by albumin contributes to a more homogeneous distribution of the drug in the body, reduces the volume of distribution, lowers the rate of clearance, and increases the plasma half-life of the drug. See ASPIRIN; BILIRUBIN; HETEROCYCLIC COMPOUNDS; THYROXINE. James A. Hamilton Bibliography. J. R. Brown and P. Shockley, Serum albumin: Structure and characterization of its ligand binding sites, Wiley, New York, 1982; J. K. Choi et al., Interactions of very long-chain saturated fatty acids with serum albumin, J. Lipid Res., 43:1000–1010, 2002; S. Curry et al., Crystal structure of human serum albumin complexed with fatty acid reveals an asymmetric distribution of binding sites, Nat. Struc. Biol., 5:827–835, 1998; C. Daniels, N. Noy, and D. Zakim, Rates of hydration of fatty acids bound to unilamellar vesicles of phosphatidylcholine or to albumin, Biochemistry, 24:3286–3292, 1985; J. A. Hamilton, Fatty acid transport: Difficult or easy, J. Lipid. Res., 39:467–481, 1998; X. M. He and D. C. Carter, Atomic structure and chemistry of human serum albumin, Nature, 358:209–215, 1992 [published erratum appears in Nature, 364(6435):362, July 22, 1993]; T. Peters, Jr., All About Albumin, Academic Press, San Diego, 1996; A. A. Spector, Fatty acid binding to plasma albumin, J. Lipid Res. 16:165–179, 1995.
Alcohol A member of a class of organic compounds composed of carbon, hydrogen, and oxygen. They can be considered as hydroxyl derivatives of hydrocarbons produced by the replacement of one or more hydrogens by one or more hydroxyl ( OH) groups. Classification. Alcohols may be mono−, di−, tri−, or polyhydric, depending upon the number of hydroxyl groups they possess. They are classified as primary (RCH2OH), secondary (R2CHOH), or ter-
tiary (R3COH), depending on the number of hydrogen atoms attached to the carbon atom bearing the hydroxyl group. Alcohols can also be characterized by the molecular configuration of the hydrocarbon portion (aliphatic, cyclic, heterocyclic, or unsaturated). There are two systems in use for alcohol nomenclature, the common naming system and the IUPAC (International Union of Pure and Applied Chemistry) naming system. The common name is sometimes associated with the natural source of the alcohol or with the hydrocarbon portion (for example, methyl alcohol, ethyl alcohol). The IUPAC method is a systematic procedure with agreed-upon rules. The name of the alcohol is derived from the parent hydrocarbon which corresponds to the longest carbon chain in the alcohol. The final “e” in the hydrocarbon name is dropped and replaced with “ol”; and a number before the name indicates the position of the hydroxyl. Examples of these two systems are given in the table. Reactions. Oxidation of primary alcohols produces aldehydes (RCHO) and carboxylic acids (RCO2H); oxidation of secondary alcohols yields ketones (RCOR). Dehydration of alcohols produces alkenes and ethers (ROR). Reaction of alcohols with carboxylic acids results in the formation of esters (ROCOR), a reaction of great industrial importance. The hydroxyl group of an alcohol is readily replaced by halogens or pseudohalogens. See ALDEHYDE; ALKENE; CARBOXYLIC ACID; ESTER; ETHER; KETONE. Uses. Industrially, the monohydric aliphatic alcohols are classified according to their uses as the lower alcohols (1–5 carbon atoms), the plasticizer-range alcohols (6–11 carbon atoms), and the detergent-range alcohols (12 or more carbon atoms). The lower alcohols are employed as solvents, extractants, and antifreezes. Esters of the lower alcohols are employed extensively as solvents for lacquers, paints, varnishes, inks, and adhesives. The plasticizer-range alcohols find their primary use in the form of esters as plasticizers and also as lubricants in high-speed applications such as jet engines. The detergent-range alcohols are used in the form of sulfate and ethoxysulfate esters in detergents and surfactants. Ethanol and methanol are volatile, combustible substances, but their use as motor fuel is limited by economic considerations. In the United States, ethanol made by fermentation of agricultural products is blended with gasoline to produce a motor fuel called gasohol, but this constitutes an insignificant fraction of the energy supply. In Brazil, where sugarcane is a major agricultural product, the large-scale production of ethanol by fermentation has been promoted by the government, and ethanol is a significant and rapidly growing component of the automotive fuel supply. See ETHYL ALCOHOL; METHANOL. Sources. Alcohols are derived either from naturalproduct processing, such as the fermentation of carbohydrates and the reductive cleavage of natural fats and oils, or by chemical synthesis based on the hydrocarbons derived from petroleum or the synthesis gas from coal.
Alcohol Alcohols and their formulas Name Common Methyl alcohol Ethyl alcohol n-Propyl alcohol Isopropyl alcohol n-Butyl alcohol sec-Butyl alcohol tert-Butyl alcohol Isobutyl alcohol n-Amyl alcohol n-Hexyl alcohol Allyl alcohol Crotyl alcohol Ethylene glycol Propylene glycol Trimethylene glycol Glycerol
IUPAC
Formula
Methanol Ethanol 1-Propanol 2-Propanol 1-Butanol 2-Butanol 2-Methyl-2-propanol 2-Methyl-1-propanol 1-Pentanol 1-Hexanol 2-Propen-1-ol 2-Buten-1-ol 1,2-Ethanediol 1,2-Propanediol 1,3-Propanediol 1,2,3-Propanetriol
CH3 OH CH3 CH2 OH CH3 CH2 CH2 OH (CH3 )2 CHO CH3 (CH2 )2 CH2 OH CH3 CH2 CHOHCH3 (CH3 )3 COH (CH3 )2 CHCH2 OH CH3 (CH2 )3 CH2 OH CH3 (CH2 )4 CH2 OH CH2 CHCH2 OH CH3 CH CHCH2 OH HOCH2 CH2 OH CH3 CHOHCH2 OH HOCH2 CH3 CH2 OH CH2 OHCHOHCH2 OH
The fermentation of sugars and starches (carbohydrates) to produce alcoholic beverages has been employed at least since history has been recorded. The industrial fermentation process is the biological transformation of a carbohydrate by a highly specialized strain of yeast to produce the desired product, such as ethanol, or 1-butanol and acetone. Fermentation is no longer the major source of 1-butanol, but still accounts for all potable ethanol and a large proportion of the ethanol used industrially worldwide. As sources of hydrocarbons based on petroleum continue to be depleted, fermentation processes based on renewable raw materials are likely to become more important. See DISTILLED SPIRITS; FERMENTATION; INDUSTRIAL MICROBIOLOGY. Syntheses. The major processes employed industrially to produce aliphatic monohydric alcohols are the reduction of synthesis gas, the hydration and oxonation of hydrocarbons, the condensation of aldehydes followed by reduction, the Ziegler process, and the reduction of animal fats and vegetable oils. Reduction. Methanol, the simplest alcohol, is produced almost entirely by the catalytic reduction of synthesis gas (carbon monoxide and hydrogen), as shown in reaction (1). +
CO Carbon monoxide
2H2
CH3OH
Hydrogen
Methanol
Ethylene CH3 CH
CH3 CH2 OH
CH3 CHCH3 OH
Propylene
(2)
Ethanol
CH2 + H2 O
Water
2-Propanol
CH2
+ CO
Ethylene
(3)
+
H2
Carbon Hydrogen monoxide CH3CH2CHO Propionaldehyde
H2
CH3CH2CH2OH (4) 1-Propanol
See FISCHER-TROPSCH PROCESS. Condensation. Some of the higher aliphatic alcohols, such as 2-ethylhexanol, are synthesized by the condensation of aldehydes followed by dehydration and reduction. This base-catalyzed condensation, known as the aldol reaction, is shown in reaction (5). OH 2CH3CH2CHO
carbons with water in the presence of an acidic catalyst is employed commercially to produce ethanol and 2-propanol, as shown in reactions (2) and (3). CH2 + H2 O
CH2
(1)
Hydration. Treatment of lower unsaturated hydro-
CH2
Most of the synthetic ethanol in the United States is made by this method. Oxonation. A modification of the Fischer-Tropsch process, known as the oxo reaction, involves the treatment of unsaturated hydrocarbons with synthesis gas at high temperatures under high pressure in the presence of a cobalt catalyst; it is a major source of monohydric aliphatic alcohols ranging from 1-propanol to the tridecanols. A mixture of isomeric alcohols is generally obtained by this process, shown in reaction (4).
base
n -Butyraldehyde
(1)
CH3CH2CH2CHCHCHO
H2O
(2) H2
CH2CH3 Butyraldol CH3CH2CH2CH2CHCH2OH
(5)
CH2CH2 2-Ethylhexanol Ziegler process. Formerly, the major source of detergent range alcohols was hydrogenolysis of animal fats and vegetable oils, but the Ziegler process, using aluminum alkyls as the catalyst and ethylene from the cracking of petroleum as the feedstock, has become the major source of industrial supply. In this sequence of reactions, a trialkyl aluminum such as
411
412
Alcohol fuel triethyl aluminum [Al(C2H5)3] is alkylated with ethylene to give the desired higher trialkyl aluminum [AlR3], which is then oxidized to the trialkoxy aluminum [Al(OR)3], which is hydrolyzed to give a mixture of primary alcohols in the desired molecular weight range. See STEREOSPECIFIC CATALYST. Ethylene glycol and propylene glycol are the most industrially significant polyhydric alcohols. They are made from ethylene, as shown in reaction (6), and O H2C
CH2
O2
Ethylene
H2C
CH2
Ethylene oxide
H2O
HOCH2CH2OH Ethylene glycol
(6) propylene, by means of their intermediate oxides. The oxides are synthesized by catalytic oxidation of the olefins. See ESTER; PHENOL; POLYOL. Paul E. Fanta Bibliography. Kirk-Othmer Encyclopedia of Chemical Technology, vol. 1, 4th ed., 1991.
Alcohol fuel Any alcohol burned as a fuel. Generally, the term alcohol refers to ethanol, the first organic chemical produced by humans. If available, any alcohol may be used as a fuel, but ethanol and methanol are the only ones which are sufficiently inexpensive. Alcohols are useful fuels, even though the oxygen atom in any alcohol molecule reduces its heating value, because the molecular structure increases the combustion efficiency. See COMBUSTION; GASOLINE. Ethanol. Alcohol (ethanol) was used as a motor fuel, either alone or in mixtures with gasoline (gasohol), in many countries during the early 1930s, and it became an important alternative fuel for automobiles again in the 1970s and 1980s. It burns well in engines designed for gasoline and has a high octane rating. However, because of the high cost of the raw materials required, or of their conversion, it costs much more than gasoline. Gasohol is produced by dissolving 5–15% absolute water-free alcohol in gasoline. The 5% water fraction in 95% alcohol causes phase separation because of its insolubility, making 95% alcohol unsuitable for use in gasohol. This is first seen as cloudiness, then as small droplets of a wateralcohol solution, and finally as two distinct layers. Such phase separation also occurs because alcohol is hygroscopic; that is, it tends to absorb water from the air and cause dilution, either before or after mixing with gasoline. Gasoline, alcohol, and water have individual solubility effects on the metals, coatings, and plastic parts of a car’s carburetor and other parts of its fuel system. Usually, gasohol has about the same effects on these materials as gasoline; and in the absence of absorbed water, gasohol causes no difficulties in usual servicing and operation of automobiles. Because alcohol increases the octane rating, and gives other improvements in combustion, gasohol usually per-
forms better and gives an increase in octane number higher than would be expected on the basis of the small fraction of alcohol in the fuel. See OCTANE NUMBER. Specially modified engines and parts in contact with water-free (absolute or neat) alcohol account also for the considerably greater efficiency obtained in its use. Neat alcohol is not used to any extent as a fuel in the United States, but is widely used as a 95% impure grade in Brazil in cars fitted with special parts. Ethanol has many advantages, but it does not compete in most countries with other liquid fuels on an economic basis for two reasons. First, raw materials are much more valuable for food or other uses, if there are no political constraints or artificial pricing. Second, basically, alcohol always requires more energy to produce than it will deliver in any use. Thus it will always be inefficient and expensive from a thermal standpoint. Nevertheless, alcohol will find use as fuel for a number of reasons: (1) In many countries there may be no other indigenous raw material from which to make motor fuels. (2) The plants for its product can be quite simple, small, and inexpensive compared to the very large, highly capital-intensive plants required for making other fluid fuels. (3) Technology is simple, well known, and established worldwide. (4) Highly labor-intensive programs are required with simple and familiar rural duties for most of the labor force. See ETHYL ALCOHOL. Methanol. Carbon in any material can be converted to methanol. Natural gas during the 1970s and 1980s was favored as a raw material, and its energy may be delivered at a much lower cost if the gas is converted to methanol than if it is piped directly for 4000–5000 mi (6400–8000 km), even less and with much greater safety than if the gas is converted to liquefied natural gas, shipped for 2000–3000 mi (3200–4800 km), and stored at its cryogenic temperature. See LIQUEFIED NATURAL GAS (LNG); NATURAL GAS. As gas prices increase, otherwise unsalable solid fuels will be used increasingly as feedstocks for methanol. Low-grade coal, lignite, or peat may be too far from market or have too much sulfur, chemically bound water, or ash to be viable fuels. Plants built alongside deposits of such fossil fuels may pipe or ship liquid methanol to market. Methanol has been synthesized in the United States since 1925, first from synthesis gas (also known as syngas, a mixture of CO and H2) manufactured from coal and, more recently, almost entirely from natural gas. While the basic chemistry is still the same, the many improvements in processing, flow sheets, catalysts, and construction materials now allow single units to produce up to 1.5 × 106 gal/day (5.7 × 106 liters/day) of fuel-grade methanol. Such so-called world-size plants are necessary if methanol is to be economical for fuel use. One such unit using, for example, low-grade lignite, will fuel automobiles for as many miles as the gasoline from 55,000 barrels (8700 m3) of crude oil a day. See LIGNITE.
Alcoholism In producing synthesis gas, the feedstock (for example, the fossil fuel), after any necessary preparation, is gasified to hydrogen, carbon monoxide, and carbon dioxide. The gasifier should operate at as high a pressure as possible to minimize the cost of compression of the synthesis gas produced up to the operating conditions of the converter: 100– 300 atm (1–3 × 107 pascals) pressure and 390–750◦F (200–400◦C). Correspondingly large is the air separation plant to supply 5 tons (4.5 metric tons) per minute of oxygen for the partial oxidation of a coal feedstock, and to give the oxygen in the large volume of methanol produced. The steps of air separation, syngas generation, balancing of reaction to produce the desired amount of hydrogen, and the separation of CO2 and H2S (then sulfur) from the synthesis gas require a vast complex of massive equipment. Purified synthesis gas passes through catalyst chambers and heat exchangers which, by correct interrelation and proper control of temperatures and space velocities, give 97% methanol, 1% water, and 1–2% higher alcohols. Modern processing (Wentworth) delivers this fuel-grade methanol without any refining being necessary for immediate use as motor fuel and with a higher energy content of the original fuel and a much lower plant cost and operating cost than earlier processes used to make the chemical grade. Methanol for fuel use has a heating value of about 10,000 Btu/lb (23 megajoules/kg). Fifteen gallons (0.57 m3) of methanol give 1 × 106 Btu (1 × 109 joules). To make it requires an input of from 1.5 to 2 times as much thermal energy based on the higher heating value of a solid fuel, or somewhat less if it is made from petroleum gas or liquid feedstocks. Thus the overall energy conversion from solid fuels is in the range of 50 to 65%, and the lower efficiencies in this range are for poorer coals or lignites. Methanol is the most versatile and cheapest liquid fuel which can be made. It is also less flammable than gasoline; accidental fires are extinguished with water instead of being spread as flaming films. When it combusts, it yields neither particulates (soot) nor sulfur oxides, and yields lower quantities of nitrogen oxides than any other fuel. In an accident at sea, there is no fire or oil slick because the methanol dissolves in water. When produced from natural gas or solid fuel, it can always be regasified to give substitute natural gas (SNG) with only a small loss of energy. Diesel engines in trucks, railroad locomotives, and ships may use a small amount of their regular diesel fuel for start-up. Then methanol may be used for at least 90% of the total fuel. Steam may be generated in almost any furnace which can be fired by oil or gas, by changing the nozzle which feeds fuel to the combustion chamber so as to use methanol, which gives no emission problems. However, methanol’s greatest potential use in electric power plants is by internal combustion in gas turbines, where it is a superior fuel. It may be used most efficiently in a combined cycle with a steam boiler and turbine on the low-temperature side. Compared with a coal-fired, steam-powered plant with its
modern accessories, methanol gives a capital cost reduced by almost 50% and a thermal efficiency over a third greater (46% versus 32%). Methanol has been preferred in automobile engines, for example, in racing cars where efficiency, performance, and safety are more important than cost. With its high octane number (110), 10% can be added to improve the performance of unleaded gasoline, along with other beneficial additives amounting to 30% or more. Methanol is used much better neat or mixed with just a few percent of common liquids as additives to eliminate problems related to cold start, materials of construction, lubricity, and so on. Lean mixtures may be burned at compression ratios of between 12 and 16 to 1, to give 50 to 80% higher thermal efficiencies than gasoline, even though the heating value of methanol is only about one-half. Even when expensive chemical-grade methanol is used in automobiles with relatively insignificant changes, the mileage costs are less than when gasoline is used as a fuel. When methanol from a world-size fuelgrade plant is used, the fuel cost per mile will be less than half. Pollution from emissions is nonexistent. See METHANOL. Donald F. Othmer Bibliography. D. F. Othmer, Chem. Eng. (London), January 1981; J. H. Perry, Methanol: Bridge to a Renewable Energy Future, 1990.
Alcoholism The continuous or excessive use of alcohol (ethanol) with associated pathologic results. Alcoholism is characterized by constant or periodic intoxication, although the pattern of consumption varies markedly. Individuals admitted for the first time to an alcoholism treatment center typically have been consuming 3–4 oz (80–100 g) of pure alcohol per day, corresponding to seven to nine drinks or bottles of beer or glasses of wine. Studies have shown that problem drinking in these populations starts at about 2 oz/day (60 g/day), that is, four to five drinks per day, and that these are consumed in rapid succession, leading to intoxication on three or more days per week. Individuals who consume these levels of alcohol have a greater-than-average risk of developing alcoholic liver cirrhosis. However, the levels should not be taken as absolute, since they can vary greatly in different individuals, according to body weight and other factors. The symptoms and consequences associated with severe alcohol consumption also vary greatly; that is, in some individuals only a few may be present. These may consist of the development of physical dependence manifested as a state of physical discomfort or hyperexcitability (tremors or shakes) that is reduced by continued consumption; the development of tolerance to the effects of alcohol, which leads individuals to increase their consumption; accidents while intoxicated; blackouts, characterized by loss of memory of events while intoxicated; work problems, including dismissal; loss of friends and
413
414
Alcoholism family association; marital problems, including divorce; financial losses, including bankruptcy or continual unemployment. Medical problems can include gastric ulcers, pancreatitis, liver disease, and brain atrophy. The last is often associated with cognitive deficiencies, as shown by the inability to comprehend relatively simple instructions or to memorize a series of numbers. See COGNITION. Individuals seeking an early treatment for their alcohol problems have very good probabilities of recovery. The lesser the number of presenting problems described above, the better the chances of favorable outcome, and so an early identification of problem drinking by family, friends, employers, or physicians becomes very important. Employee assistance programs have become an important factor in identification and referral and rehabilitation of individuals with alcohol problems in the United States, Canada, and many other countries. The types of intervention vary greatly, progressing from self-monitoring techniques, to intensive outpatient and inpatient programs, to Alcoholics Anonymous groups. Absorption of alcohol. Alcohol is absorbed most rapidly from solutions of 15–30% (30–60 proof) and less rapidly from beverages containing below 10% and over 30%. This is to be expected at the lower concentrations, since the rate of absorption depends on the concentration gradient across the mucosal surface. At higher concentrations, ethanol abolishes the rhythmic opening of the pylorus, thus preventing the passage to the intestine, where absorption is faster. The presence of food in the stomach is also known to delay gastric emptying and thus to slow absorption. Once in the bloodstream, alcohol is distributed evenly in all tissues, according to their water content. As a rule of thumb, for a 150-lb (68-kg) individual two standard drinks (1.5 oz of a distilled beverage or 13.6 g per standard drink) will yield blood alcohol levels of about 0.06 g per 100 mL of blood. It was previously believed that all ethanol consumed was absorbed into the bloodstream. However, studies have shown that a small part of the ethanol ingested is degraded directly in the stomach without entering the blood. See DISTILLED SPIRITS; MALT BEVERAGE; WINE. Effects on the nervous system. The exact mechanisms of the pharmacological actions of alcohol are not known. Alcohol can act as a stimulant at lower doses, and as a depressant at higher doses. Even at very low doses, alcohol can impair the sensitivity to odors and taste. Also, low doses are known to alter motor coordination and time and space perception, important aspects of car driving (about 50% of all fatal traffic accidents are caused by intoxicated drivers). Some effects are already seen at levels of 0.05%. Pain sensitivity is diminished with moderate doses. In some individuals, alcohol is known to diminish feelings of self-criticism and to inhibit fear and anxiety, effects which are probably related to an alcohol-induced sociability. These effects act, no doubt, as psychological reinforcers for the use of alcoholic beverages.
It is generally accepted that alcohol affects the nerve cell by preventing the production and propagation of electric impulses along a network consisting of axons and synapses. The brain functions much as an electronic system in which one nerve cell, acting as a current generator, communicates information to many other cells, which in turn receive impulses from many other areas. Some impulses are enhanced, others are blunted. Memory and conditioning appear to play an important role in integrating the impulses which are finally expressed as behaviors. Studies in the United States and England have shown that when alcohol becomes dissolved in the membrane of the cells, it fluidizes or disorganizes the membrane, which in turn leads to changes in the physical and biochemical characteristics of the latter. Chronic exposure to alcohol alters the composition of the membrane and its rigidity, so that alcohol becomes less of a disorganizing agent. The new membrane composition also appears to partly exclude the alcohol molecules from dissolving in it. These changes are likely to be related to tolerance, when more alcohol is required to produce the same effect. Studies have shown that learning and conditioning also play a role in the latter phenomenon. In general, the knowledge gained from these studies is that the state of organization (fluidity–rigidity) of the membrane is correlated with the state of excitability of the brain; the increased rigidity of the membrane that follows chronic alcohol consumption might account for the alcohol withdrawal hyperexcitability (physical dependence) and the fact that continued ingestion of alcohol is required to alleviate the physical dependence. Studies in rodents and monkeys have addressed the question of what factors can lead to alcohol consumption in these animals, which normally reject alcohol. It has been found that small amounts of alcohol-derived molecules, such as the tetrahydroisoquinolines, when injected into the brain, produce a marked increase in the preference for alcohol of the animals. These studies raise the possibility that these substances could be formed in greater amounts in individuals who consume alcohol heavily. Studies have shown that lipid molecules in the membranes of brain cells of animals fed alcohol for a prolonged time can confer alcohol resistance to brain cell membranes of animals that have never received alcohol. Not only brain lipids are modified but also some brain proteins. Thus, it is clear that brain chemistry is modified by chronic alcohol exposure. This effect remains for days or weeks after alcohol administration has been discontinued. Brain cells communicate with each other and with cells involved in sensation by means of chemical messengers, or neurotransmitters, released by these cells. Once a neurotransmitter is externally recognized by a cell, the information is passed to the interior of the cell by second messengers, a system that involves G-proteins. The continued presence of alcohol makes the nerve cells produce less of these proteins, and as a consequence an external chemical message is not well communicated to the interior of the cell.
Alcoholism Studies using white blood cells of alcoholics as sensors of this deficiency have shown that as much as 75% of the external message is not transmitted to the interior of the cell. These changes are of importance in that they may explain the reduction in cognitive efficiency, in relation to learning and to processing external information, which starts at levels of alcohol consumption that are considered moderate social drinking. A major finding in the mid-1980s was that some of the effects of alcohol can be quickly reversed by new experimental drugs. Studies have shown that alcohol enhances the actions of an inhibitory brain neurotransmitter referred to as gamma-aminobutyric acid (GABA). Benzodiazepines, such as diazepam, are anxiety-reducing and sedative drugs which also enhance the effects of GABA. These effects can be reduced by experimental antagonist molecules, which interact in the brain in the same regions where GABA is found. See CELL MEMBRANES; NERVOUS SYSTEM (VERTEBRATE); SYNAPTIC TRANSMISSION. Effects on the liver. The liver is responsible for about 80% of the metabolism of alcohol. In the liver, alcohol is first oxidized to acetaldehyde and then to acetate, which is metabolized in many tissues, including the brain, heart, and muscles. A 150-lb (68-kg) person metabolizes approximately 0.4 oz (10 g) of pure alcohol per hour (about 1 oz of a distilled beverage per hour) or, if alcohol is continuously present in the bloodstream, about 8–10 oz (190–240 g) of pure alcohol per day, equivalent to 1300–1600 calories per day. Since alcoholic beverages contain negligible levels of essential nutrients, these calories are called empty calories. Many alcoholics show malnutrition due to the fact that an important part of their caloric intake is alcohol. Alcohol also impairs the absorption and the metabolism of some essential nutrients. See ABSORPTION (BIOLOGY); LIVER; MALNUTRITION. In the presence of alcohol, about 80% of oxygen consumed by the liver is devoted to the metabolism of alcohol; as a consequence, other substances such as fats, normally oxidized by the liver, are not metabolized, leading to fat accumulation in the liver. Continued consumption of alcohol has been reported to increase the capacity of the liver to oxidize alcohol, such that some alcoholics can metabolize up to 16 oz (400 g) of pure alcohol per day. Since all the hepatic metabolism of alcohol, by any known route, requires oxygen, there is also an increased uptake of oxygen by the liver. See LIPID METABOLISM. Alcoholic liver disease and liver cirrhosis rank among the 10 leading causes of mortality in the United States and Canada. It was initially believed that cirrhosis in the alcoholic was exclusively of nutritional origin. However, baboons fed 50% of their caloric intake as alcohol develop cirrhosis in 1–4 years, even when the alcohol is given with a nutritious diet. Baboons fed the same diet, but with carbohydrates given instead of alcohol, did not develop cirrhosis. It is not known, however, if nutrient supplementation would have reduced the damage in its initial stages. Some investigators have proposed that in humans alcoholic liver disease is
produced by increased utilization of oxygen by the liver, when the extra utilization is not matched by an adequate supply of oxygen to the organ. Studies in Canada have shown that the antithyroid drug propylthiouracil, which reduces the increased oxygen demand induced by alcohol in the liver, greatly protects against the mortality of alcoholic liver disease. In a 2-year treatment study, deaths in patients treated with propylthiouracil were 50–60% lower than in patients given placebo capsules. The mortality due to liver disease can be drastically reduced if the individuals abstain from alcohol. Since it is unlikely for liver scar tissue in cirrhosis to reverse to normal, studies suggest that complications accompanying cirrhosis may be more damaging than the cirrhosis itself. Alcoholic liver disease is characterized by two conditions: failure of the liver to detoxify noxious substances and to produce essential products, and increased resistance to blood flow through the liver. The latter condition (portal hypertension) results in the opening of collateral vessels (esophageal varices) as an alternative for the passage of portal blood. The expanded vessels can burst, leading to internal bleeding. Also, noxious substances bypass the cleansing action of the liver. This condition is responsible for about one-half of the deaths in alcoholic liver disease. Studies have demonstrated that chronic alcohol consumption results in enlargement of the liver cells, which leads to compression of the minute intrahepatic vessels that feed the liver cells. In turn, compression of these vessels increases the resistance to blood passage and to portal hypertension. This effect is quickly reversible upon abstinence from alcohol, and may be responsible for the reduction in mortality from liver disease in abstinent alcoholics. See CIRRHOSIS. Effects on sexual and reproductive functions. Alcohol consumption produces a striking reduction in the hormone testosterone in males. This appears to be due to an effect in brain areas which are responsible for specific signals for the maintenance of normal testosterone levels and to a metabolic effect in the testicular cells that produce testosterone. In addition, alcoholics have increased rates of testosterone destruction by the liver. Feminization signs are often seen in male alcoholics, while impotence and testicular atrophy are also found frequently. Studies in France and the United States in the early 1970s were the first to report a distinct pattern of physical and behavioral anomalies in children of alcoholic women who drank heavily during pregnancies. These anomalies which characterize the fetal alcohol syndrome include facial and cranial deformities such as drooping eyelids (ptosis), small eyes (microophthalmia), and underdeveloped midface (midfacial hypoplasia). A 10-year follow-up of these children in the 1980s indicated that all were borderline for intelligence, while half of them presented IQs of 20 to 60, in the retarded range. The prevalence of this condition in the general population has not been fully established, but it ranges between 0.5 and 3 births in 1000. In the United States the risk of fetal alcohol
415
416
Alcoholism syndrome is about seven times higher in the black race than in the white population. Studies so far do not indicate that there is a safe level of habitual alcohol consumption, below which there is no effect on the unborn; and the role of specific drinking patterns, including a single bout, needs to be determined. In animals with short gestation periods, a single massive intoxication results in damage to the newborn. Social drinking during pregnancy has been associated with lower scores on tests of spoken English and verbal comprehension in the children. See FETAL ALCOHOL SYNDROME. Early identification of alcohol abuse. Several laboratory tests are used to detect high levels of alcohol consumption. Among these are an elevated serum gamma-glutamyl transferase and the presence of enlarged red cells. In addition, rib and vertebral fractures, as seen on routine x-rays, are found 15 times more frequently in alcoholics than in normal populations. Combinations of these tests allow diagnosis of heavy alcohol consumption with a good degree of accuracy. A simple litmus-paper-like test to detect alcohol in saliva, serum, and urine consists of a strip of paper with special reagents which, when impregnated by any of the fluids, indicates the amount of alcohol present by changing color. Another test to determine intoxication involves the eyealyzer, which in 15 seconds senses alcohol vapors emanating from the lacrimal fluid normally covering the eye surface. This noninvasive device is 95% accurate and is not affected by residual alcohol in the mouth, as is the breathalyzer, and it can be used in persons that are unconscious since the eyelid can be held open. Genetic factors. There is abundant evidence that tendency to alcoholism can be of familial origin, due to environmental, cultural, and genetic factors. A Swedish study demonstrated that identical twins are twice as likely to have a common alcoholic problem as fraternal twins. In an American-Danish study, it was shown that children of alcoholic parents are more likely to develop alcoholism (18%) than children of nonalcoholic parents (4%) when both groups of children were adopted by nonrelatives within 6 weeks of birth. In a Swedish study, similar findings were observed where male adoptees whose biological fathers were severely alcoholic had a 20% occurrence of alcohol abuse, as compared with 6% for adoptee sons of nonalcoholic fathers. All adoptions in this study had occurred with the first 3 years of life and in most cases in the first few months. A subsequent Swedish-American study revealed the existence of two types of genetic predisposition to alcoholism. The more common type, mileu-limited alcoholism (type I), is associated with mild manifestations and adult-onset alcohol abuse in one of the biological parents. In this group, a significant contribution to alcoholism predisposition was low socioeconomic status. A second type, male-limited alcoholism (type II), accounts for about 25% of all male alcoholics and is unaffected by the environment. This type of alcoholism, often with teenage onset, was associated with severe alcoholism in the biological father but
not in the mother. The biological father also tended to have a record of more serious criminality and an earlier onset of alcoholism. Alcohol abuse was nine times more frequent in adoptees of type II alcoholic fathers, regardless of the environment into which they were adopted. Brain electrical activity in response to external stimuli is markedly different in the male children of type II alcoholics, at an age as early as 12 years. It should be noted that about 60–80% of sons of alcoholics do not become alcoholic and thus environmental factors have a marked influence, for example in type I alcoholism. In identical twins, where genetic factors are constant, alcoholism occurs in both individuals in only 25% of the cases. Studies from Germany and Japan demonstrate that some genetic factors can protect against the development of alcoholism. About one-half of the population of Japan lack a specific aldehyde dehydrogenase to rapidly metabolize the toxic acetaldehyde produced in the liver in the first reaction of alcohol oxidation. As a consequence, blood acetaldehyde levels build up, leading to facial flushing and to unpleasant reactions which may include nausea. These individuals appear not to develop alcoholism, as shown by the finding that an extremely low proportion of individuals lacking this specific aldehyde dehydrogenase seek treatment of alcoholism. A number of studies also showed that special genetic strains of rodents can be developed which either prefer alcohol to water or reject alcohol. Some strains of alcohol-preferring rats drink to intoxication, develop physical dependence manifested as hyperexcitability upon withdrawal from alcohol, and have been shown to work (conduct specific tasks) to receive alcohol. Other genetic strains of mice can be shown to differ markedly in sensitivity to the effects of ethanol in that after receiving similar doses of ethanol, one strain (LS or long-sleep) sleeps five times longer than another strain (SS or short-sleep). These differences have been found to result from differences in the brain of these animals. See BEHAVIOR GENETICS. Statistics. In the United States there are an estimated 10 million problem drinkers. Official statistics indicate that in 1986 1.2 million people were admitted for alcoholism treatment in private and government facilities. Alcohol-related deaths are estimated at 100,000 per year. Accidents and cirrhosis are the leading causes of mortality associated with excessive alcohol use, followed by hypertensive diseases, infections, and cancer. Eighty percent of all cirrhosis and about one-half of all accidental deaths, suicides, and homicides are alcohol-related. The average per-capita consumption (over 14 years of age) is 2.65 gallons (10 liters) of pure alcohol per year. Heavy drinkers, constituting about 10% of the drinking population, account for one-half of all alcohol consumed in the United States. It is estimated that alcohol creates problems for 18 million persons 18 years old or more. Both the level of consumption and the number of deaths attributed to alcohol declined during the 1980s. Women drink significantly less than men; however, there has been
Alcoholism an increase in drinking among women aged 35 to 64. People over age 65 consume less alcohol than younger adults and have lower prevalence of alcohol abuse. Studies in the Scandinavian countries and Canada have shown that the price of alcoholic beverages (relative to the disposable income) is inversely related to the total amount of alcohol consumed by the population. Cirrhosis in different countries has been shown to correlate almost perfectly with the percapita consumption, with the United States showing about one-third of the incidence found in France. Some studies have shown that the frequency of coronary heart disease is lower in countries with higher per-capita consumption of alcohol. However, others suggest that this is largely due to the fact that populations in these countries also have dietary patterns that by themselves lead to less coronary heart disease. See ADDICTIVE DISORDERS; BEHAVIORAL TOXIYedy Israel COLOGY. Pharmacotherapy. Pharmacotherapy for alcohol rehabilitation has been gaining wider acceptance. Specific pharmacotherapies which have received the most research attention utilize naltrexone, acamprosate, and disulfiram. Other pharmacological interventions for specific indications are SSRIs, (selective serotonin reuptake inhibitors) and buspirone. Naltrexone. Naltrexone is an opiate receptor antagonist which blocks the effects of endogenous opioids in the brain. Research from animal studies suggests that alcohol activates endogenous opioid systems and may contribute to the pleasurable effects produced by alcohol consumption. Consequently, naltrexone might reduce the reinforcing effects of alcohol consumed by people and decrease their incentive to drink. The efficacy of naltrexone treatment has been examined by several investigators in double-blind, placebo-controlled studies. These studies have demonstrated that naltrexone (50 mg per day) reduces the frequency of alcohol consumption, resumption of heavy drinking, and intensity of alcohol craving with relatively few significant side effects. Relative to placebo controls, relapse rates were reduced significantly when patients took naltrexone. Naltrexone is well tolerated; about 15% of patients suffer from nonserious side effects, primarily nausea. These findings contributed to the 1995 decision by the U.S. Food and Drug Administration (FDA) to approve naltrexone for use in alcoholism treatment. Acamprosate. The mechanism of action of acamprosate, calcium-acetyl-homotaurinate, involves primarily the restoration of a normal N-methyl-Daspartate (NMDA) receptor tone in glutamatergic systems. It decreases postsynaptic potentials in the neocortex and diminishes voluntary alcohol intake in alcohol-preferring rats. It was proposed that one of acamprosate’s actions is suppressing conditioned withdrawal craving by its effects on NMDA receptors and calcium channels. Acamprosate was investigated in 20 controlled clinical trials with about 5000 patients up to early 2006. Acamprosate (1998 mg per day) lengthens time to relapse, reduces drinking days, and increases complete abstinence
among alcohol-dependent patients. Adverse events tended to be mild and transient, primarily involving the gastrointestinal tract with diarrhea and abdominal discomfort in about 10% of the patients. In 2004, acamprosate was approved in the United States. Combined treatment. It was expected that naltrexone would specifically attenuate “reward craving” (anticipating hedonic effects of alcohol), while acamprosate would diminish relief craving (anticipating unpleasant effects associated with the absence of alcohol). The combination of these two drugs would act simultaneously on two different aspects of craving and might have a more incisive effect on the risk of relapse than either treatment alone. Both drugs are well tolerated and have no marked propensity for potentially dangerous drug–drug interactions. Now preclinical and clinical data exist to demonstrate a good tolerability of combined medication and, moreover, provide some hints for a significant improvement in outcome regarding the duration of abstinence following alcohol withdrawal. Disulfiram. Disulfiram is a drug which causes an inhibition of the enzyme aldehyde dehydrogenase. Aldehyde dehydrogenase breaks down the first metabolite of alcohol, acetaldehyde. After taking disulfiram for 3 or more days, a person who drinks alcohol usually will experience an increase in acetaldehyde blood levels. This rise will produce nausea, vomiting, tachycardia, difficulty in breathing, and changes in blood pressure leading to hypotension. These effects typically last up to 3 h. When the drug is prescribed at the standard dose (250 mg per day), dangerous side effects are minimized, although the side-effect profile of disulfiram must be considered carefully for each patient. Given its pharmacological effects, disulfiram acts as a deterrent to future alcohol consumption by making intake an aversive experience. Alcohol use cannot be resumed safely until 6–7 days after the last disulfiram dose. Despite the compelling rationale for its potential effectiveness, studies of disulfiram’s efficacy are equivocal, providing only modest evidence for reducing drinking frequency without significantly enhancing abstinence rates. While some professionals endorse the effectiveness of disulfiram treatment, others advocate more controlled clinical trials to determine its relative costs and benefits before it is used routinely in alcoholism treatment. Several new drug therapies for alcoholism rehabilitation are under investigation. Buspirone, a nonbenzodiazepine antianxiety agent, may decrease anxiety symptoms associated with a protracted alcohol withdrawal syndrome, thus reducing alcohol relapse potential. SSRIs also seem to reduce alcohol intake, especially in clinically depressed patients suffering from a co-morbid alcohol dependence. These medications require further investigation to determine their effectiveness as pharmacotherapeutic agent, in the treatment of alcoholism. Pharmacotherapy approaches for alcohol rehabilitation are beginning to show promise as clinical interventions assisting patients in their efforts to recover from alcoholism. Concurrent behavioral and
417
418
Alcyonacea supportive psychotherapeutic treatments, self-help program participation, and supervision of prescribed medication regimens may further optimize pharmacotherapy outcomes. Determining the best combination of pharmacotherapy approaches, psychotherapeutic interventions, and patient characteristics remains a challenge to the alcoholism treatment field. Falk Kiefer; Steve Martino Bibliography. H. Bogleites and B. Kissin (eds.), The Genetics of Alcoholism, 1994; Combine Study Group, Testing combined pharmacotherapies and behavioral interventions for alcohol dependence (the COMBINE study): A pilot feasibility study, Alcohol Clin. Exp. Res., 27(7):1123–1131, 2003; J. C. Garbutt et al., Pharmacological treatment of alcohol dependence: A review of the evidence, JAMA, 281:1318–1325, 1999; D. B. Goldstein, Pharmacology of Alcohol, 1983; D. W. Goodwin, Alcoholism, 2d ed., 1994; J. C. Hughes and C. C. Cook, The efficacy of disulfiram: A review of outcome studies, Addiction, 92:381–395, 1997; F. Kiefer et al., Comparing and combining naltrexone and acamprosate in relapse prevention of alcoholism: A double-blind, placebo-controlled study, Arch. Gen. Psychiat., 60:92–99, 2003; R. Z. Litten and J. P. Allen, Pharmacotherapies for alcoholism: Promising agents and clinical issues, Alcoholism Clin. Exp. Res., 15:620–633, 1991; C. P. Rivers, Alcohol and Human Behavior: Theory, Research, and Practice, 1994; J. R. Volpicelli et al., Effect of naltrexone on alcohol “high” in alcoholics, Amer. J. Psychiat., 152:613–615, 1995.
Alcyonacea An order of the cnidarian subclass Alcyonaria (Octocorallia). Alcyonacea, the soft corals (see illus.), are littoral anthozoans, which form massive or dendriform colonies with yellowish, brown, or olive colors. Most attach to some solid substratum; however, some remain free in sandy or muddy places. The only skele-
Alcyonaceans. (a) Alcyonium palmatum. (b) Dendronephytha sp.
tal structures are small, elongated, spindle-shaped or rodlike, warty sclerites which are scattered over the mesoglea. The colony is supple and leathery. The polyp body is embedded in the coenenchyme, from which retractile anthocodia protrude in Alcyonium. In Xenia and Heteroxenia, anthocodia are nonretractile. The polyp base is protected by many sclerites and is termed a calyx. The aciniform gonads develop in the six ventral mesenteries and hang inside the gastrovascular cavity. Anthomastus, Sarcophyton, and Lobophytum are dimorphic and the siphonozooid is fertile only in Anthomastus. New polyps arise asexually also, from the solenial network. See OCTOCORALLIA (ALCYONARIA); CNIDARIA. Kenji Atoda
Aldebaran A cool red giant star, prominently located in the constellation Taurus. With its red color, it appropriately represents the eye of the Bull. At a distance of 18.5 parsecs (60 light-years), Aldebaran, or α Tauri, is among the nearest (and brightest) giant stars to the Sun. It appears to be situated near the middle of the Hyades, but Aldebaran is actually less than half the distance of this nearby star cluster. The star is an example of a K-type giant, a very common type of evolved star that derives its energy from the thermonuclear burning of helium in a core surrounded by a thin, hydrogen-burning shell. Its spectral type of K5III corresponds to an effective temperature of 6700◦F (4000 K) and a radius of about 40 times the Sun. It is nearly 150 times more luminous than the Sun and, as is typical for K giants, its brightness varies by a modest amount. Aldebaran is accompanied in a long-period binarystar system by a cool dwarf companion star some 100,000 times fainter than the giant. See BINARY STAR; GIANT STAR; HYADES; SPECTRAL TYPE; STELLAR EVOLUTION; TAURUS; VARIABLE STAR. Because it is among the nearest and brightest members of its class, Aldebaran has been studied in considerable detail. A thin and very cool shell of dust surrounds the star out to 500 stellar radii. The dust grains originate in the outer atmosphere of Aldebaran and are then gently blown away by radiation pressure. Aldebaran’s angular diameter of about 20 milli-arc-seconds (1/180,000 degree) has been measured by several interferometric methods, including a measurement with a two-telescope interferometer that determined the diameter to an accuracy of less than 1%. Observations from the Hubble Space Telescope detected the emission from ionized atoms in the chromosphere of Aldebaran, where the turbulent gases appear to be churning at speeds of more than 12 mi/s (20 km/s). See STAR. Harold A. McAlister Bibliography. A. Frankoi, D. Morrison, and S. C. Wolff, Voyages Through the Universe, 3d ed., Brooks/Cole, 2004; J. B. Kaler, The Hundred Greatest Stars, Copernicus Books, 2002; J. M. Pasachoff and A. Filippenko, Astronomy in the New
Aldehyde Millennium, 3d ed., Brooks/Cole, 2007; A. Quirrenbach et al., Angular diameter measurements of cool giant stars in strong TiO bands and in the continuum, Astrophys. J., 406:215–219, 1993.
peracids, silver oxide, and nitric acid, as in reaction (1). O
O
RCH + (O)
Aldehyde One of a class of organic chemical compounds represented by the general formula RCHO. Formaldehyde, the simplest aldehyde, has the formula HCHO, where R is hydrogen. For all other aldehydes, R is a hydrocarbon radical which may be substituted with other groups such as halogen or hydroxyl (see table). Aldehydes are closely related to ketones, since both types of compounds contain the carbonyl group C
O
In the case of ketones, the carbonyl group is attached to two hydrocarbon radicals. Because of their high chemical reactivity, aldehydes are important intermediates for the manufacture of resins, plasticizers, solvents, dyes, and pharmaceuticals. See KETONE. Physical properties. At room temperature formaldehyde is a colorless gas, which is handled commercially in the form of a 37% aqueous solution or as its polymer, paraformaldehyde (CH2O)x. The other low-molecular-weight aldehydes are colorless liquids having characteristic, somewhat acrid odors. The unsaturated aldehydes acrolein and crotonaldehyde are powerful lacrimators. Chemical properties. The important reactions of aldehydes include oxidation, reduction, aldol condensation, Cannizzaro reaction, and reactions with compounds containing nitrogen. Oxidation. Aldehydes are oxidized to carboxylic acids by a wide variety of oxidizing agents, such as air, oxygen, hydrogen peroxide, permanganate,
Aldehydes and their formulas Compound Formaldehyde Acetaldehyde Acetaldol Propionaldehyde n-Butyraldyde Isobutyraldehyde Acrolein Crotonaldehyde Chloral Chloral hydrate Benzaldehyde
(1)
The catalytic oxidation of acetaldehyde is used for the production of three industrially important compounds, as in reaction (2), where the product CH3CHO
O2 (air)
O
catalyst
CH3COOH, (CH3CO)2O, CH3COOH Acetic acid
(2)
Acetic Peracetic anhydride acid
depends on the conditions of the reaction. Reduction. Aldehydes are reduced to primary alcohols cheaply by reaction with hydrogen in the presence of a platinum or nickel catalyst, or in the laboratory by treatment with aluminum isopropoxide or oxide or with complex metal hydrides such as lithium aluminum hydride (LiAlH4) or sodium borohydride (NaBH4), as in reaction (3). RCHO + (H) −→ RCH2 OH
(3)
Aldehydes react with hydrogen in the presence of ammonia and a catalyst to form amines. This process [reaction (4)], known as reductive amination, is used Pt or Ni
RCH O + NH3 + H2 Aldehyde
RCH2 N H2 + H2 O Amine
(4)
commercially to produce many important amines. Aldol condensation. Two molecules of an aldehyde react in the presence of alkaline, acidic, or amine catalysts to form a hydroxyaldehyde. This reaction proceeds only when the carbon adjacent to the carbonyl group has at least one hydrogen atom [reaction (5)]. OH
Formula
catalyst
RCH2CHO + RCH2CHO
HCHO CH3 CHO CH3 CHOHCH2 CHO C2 H5 CHO CH3 (CH2 )2 CHO (CH3 )2 CHCHO CH2 CHCHO CH3 CH CHCHO CCI3 CHO CCI3 CH(OH)2
RCH2CHCHCHO
Aldehyde
(5)
R Aldol (a hydroxy aldehyde)
Aldols are readily dehydrated to form unsaturated aldehydes, and the latter can be hydrogenated to form saturated, primary alcohols. This reaction sequence is the basis of the commercial preparation of crotonaldehyde and 1-butanol (n-butyl alcohol) from acetaldehyde [reaction (6)].
H C
Cinnamaldehyde
RC — OH
O
H
H
H
C
C
C
2CH3CHO
alkaline catalyst
CH3CHOHCH2CHO
heat
Acetaldol O CH3CH
CH
CHO
Crotonaldehyde
H2
CH3CH2CH2OH 1-Butanol
(6)
419
420
Aldehyde Cannizzaro reaction. On treatment with strong alkali, aldehydes which lack a hydrogen attached to the carbon adjacent to the carbonyl group undergo a mutual oxidation and reduction, yielding one molecule of an alcohol and one molecule of the carboxylic acid salt. The reaction of benzaldehyde with aqueous potassium hydroxide is typical [reaction (7)].
Since these aldehyde derivatives are usually crystalline compounds with sharp melting points, they are frequently employed for the characterization and recognition of aldehydes. See HYDRAZINE; OXIME. Cyanohydrin formation. Aldehydes combine with hydrogen cyanide to form hydroxynitriles, known as cyanohydrins, as shown in reaction (14) for benz-
H O
2 2C6 H5 CHO + KOH − −−→ C6 H5 CH2 OH + C6 H5 COOK (7)
Acetaldehyde and formaldehyde in the presence of calcium hydroxide react by an aldol condensation followed by a “crossed” Cannizzaro reaction to give the industrially useful polyhydroxy compound pentaerythritol [reaction (8)]. CH3CHO + 4CH 2O + Ca(OH)2 + H2O Formaldehyde CH2OH HOCH2
C
CH2OH + (HCOO)2Ca
RCHOH
(9)
C6 H5 CHO + N2 NC6 H5 −→ NC6 H5 + H2 O (10)
Oximes, phenylhydrazones, and semicarbazones. Aldehydes react readily with hydroxylamine to eliminate water and form oximes [reaction (11)]. RCH NOH + H2 O Oxime
H2 NOH
RCH
(15)
(11)
Phenylhydrazine and semicarbazide react in a similar manner to produce phenylhydrazones and semicarbazones, respectively [reactions (12) and (13)].
Hemiacetal
Hemiacetals may react further with alcohols in the presence of acidic catalysts, but not basic catalysts, to form acetals [reaction (16)]. OCH3
OCH3 RCH
However, aromatic aldehydes such as benzaldehyde react with primary amines such as aniline with elimination of a molecule of water to give stable and well-characterized imines, also known as Schiff bases [reaction (10)].
Hydroxylamine
OCH3
Aldehyde Methanol
NH2
RCHO +
aldehyde. This particular cyanohydrin, known as mandelonitrile, occurs in nature in the form of a glycoside derivative, amygdalin, which is responsible for the characteristic flavor of bitter almonds and the seeds of peaches and apricots. Hemiacetal-acetal formation. Alcohols react with aldehydes in the presence of acidic catalysts to form hemiacetals [reaction (15)].
OH Calcium formate
Reaction with ammonia and amines. The addition of ammonia to the carbonyl group of an aldehyde is reversible, and the resulting aminoalcohol usually cannot be isolated [reaction (9)].
Aldehyde
(14)
OH
RCHO + CH3OH
Pentaerythritol
C6 H5 CH
C6H5CHCN
(8)
CH2OH
RCHO + NH3
C6H5CHO + HCN
+ CH3OH
RCH
OH
+ H2O
(16)
OCH3
Hemiacetal Methanol
Acetal
Hemiacetals are generally too unstable to permit isolation, but acetals are very stable under basic conditions, and they are therefore useful as “protecting groups,” preventing reaction of the carbonyl group while transformations are carried out elsewhere in the molecule. Condensation with active methylene compounds. Aldehydes condense with a variety of active methylene compounds in the presence of a base with elimination of a molecule of water to give unsaturated acids or acid derivatives. The reactions of benzaldehyde shown in reactions (17)–(19) are typical. C6H5CHO + CH2(COOH)2 Malonic acid C6H5CH
C(COOH)2
(17)
Benzylidenemalonic acid
RCHO + C6H5NHNH2 Aldehyde Phenylhydrazine RCH NNHC6H5 + H2O (12) Phenylhydrazone RCHO + H2NCONHNH2
C6H5CH
CHCOOH
(18)
Cinnamic acid
C6H5CHO + CH3CO2C2H5
Aldehyde Semicarbazide RCH
C6H5CHO + (CH3CO)O Acetic anhydride
Ethyl acetate NNHCONH2 + H2O
Semicarbazone
(13)
C6H5CH
CHCO2C2H5
Ethyl cinnamte
(19)
Aldehyde Wittig reaction. Aldehydes react with phosphorus ylides to form olefins, as shown in reaction (20). The
RCHO + (C6H5)3P
C
R'
R" Phosphorus ylide
Aldehyde
RCH
R' + (C6H5)3PO
C
(20)
hyde is produced similarly by the dehydrogenation of ethanol over a copper catalyst at 250– 300◦C (480–570◦F) with formation of hydrogen as a by-product. Oxidation of primary alcohols. Acetaldehyde is produced on a large scale industrially by passing a mixture of ethanol, oxygen, and steam over silver gauze at 480◦C (896◦F) to give an overall yield of 85–90% of the aldehyde [reaction (23)]. Ag
CH3 CH2 OH + 1/2 O2 −−− −−→ CH3 CHO + H2 O
R" Olefin
Phosphine oxide
aldehyde may be aliphatic or aromatic. It may contain double or triple bonds and functional groups such as hydroxyl, alkoxyl, nitro, halo, acetal, and ester. Thus the reaction is very general, and the position of the double bond is always certain. Reformatsky reaction. Aldehydes react with bromo esters and zinc to form hydroxy esters [reaction (21)]. RCHO + BrCH2COOC2H5 + Zn Aldehyde
Ethyl bromoacetate
On a laboratory scale, oxidizing agents such as chromic acid or manganese dioxide have been used to produce aldehydes from primary alcohols, as in reaction (24). The reaction must be carefully con(O)
RCH2 OH −−− −−→ RCHO
H2 C
(21)
Pd,CuCl
2 CH3 CHO CH2 + 1/2 O2 −−−−−−→
Hydroxy ester
Benzoin condensation. Aromatic aldehydes undergo
bimolecular condensation in the presence of alkali cyanide catalyst [reaction (22)].
CH2
CHCH3 + O2
Propylene
Catalyst 300–400˚C (570–750˚F)
CH2
O
Benzaldehyde
NaCN
(25)
Acrolein, prepared industrially by the oxidation of propylene over a copper oxide catalyst at 350◦C (660◦F), as in reaction (26), is used in the manufac-
CH2COOC2H5
2C6H5CHO
(24)
trolled to avoid oxidation of the aldehyde to the carboxylic acid. Oxidation of olefins. Ethylene is oxidized directly to acetaldehyde by the Wacker process, employing a palladium–cupric chloride catalyst as shown in reaction (25).
OH RCH
(23)
CHCH
O + H2O
(26)
Acrolein C6H5CHOHCC6H5
(22)
Benzoin
Tests for aldehydes. Tollens’ reagent, a solution of silver nitrate in ammoniacal sodium hydroxide, oxidizes aldehydes, with the formation of a mirror of metallic silver. Schiff’s reagent (fuchsin aldehyde reagent) reacts with aldehydes to give a violet-purple color. Fehling’s solution (a blue, alkaline cupric tartrate complex) is especially valuable for the characterization of reducing sugars. A positive test is indicated by the formation of a red precipitate of cuprous oxide. In the infrared absorption spectrum, the aldehyde group has a characteristic strong band at 1660– 1740 cm−1 due to the carbonyl stretching vibration accompanied by two weak bands at 2700–2900 cm−1 due to the aldehydic carbon-hydrogen stretching. In the proton magnetic resonance spectrum the aldehydic protons have a unique absorption region at 0.0– 0.6 τ . Synthesis. Because of the importance of aldehydes as chemical intermediates, many industrial and laboratory syntheses have been developed. The more important of these methods are illustrated by the following examples. Catalytic dehydrogenation of primary alcohols. Formaldehyde is produced on a large scale industrially by the catalytic dehydrogenation of methanol. Acetalde-
ture of glycerol and acrylic acid. See GLYCEROL. In the laboratory, olefins are treated with ozone to form an ozonide which decomposes with water to form aldehydes and hydrogen peroxide [reaction (27)]. O RCH
CHR' + O3
RCH
CHR' O
O Olefin
H2O
Ozonide RCHO + R'CHO + H2O2 (27)
Hydroformylation of olefins. Olefins may be hydroformylated by reaction with carbon monoxide and hydrogen in the presence of a catalyst, usually cobalt carbonyl, at 180◦C (360◦F) and 6000 lb/in.2 (41 megapascals) as shown in reaction (28). This reaction is
CH2
CH2 + CO + H2
Ethylene
Co2(CO)8
CH3CH2CHO
(28)
Propionaldehyde
also known as the oxo process. Chlorination of a methyl group attached to a benzene ring. The chlorination of toluene in the presence of strong light (which serves as a catalyst) proceeds stepwise
421
422
Aldehyde to form benzyl chloride and benzal chloride at 135– 175◦C (275–347◦F), by reactions (29) and (30), respectively. C6H5CH2Cl + HCI
C6H5CH3 + Cl2 Toluene
(29)
Benzyl chloride
tory method, since the acid chlorides can be prepared from the corresponding carboxylic acids. Reaction of Grignard reagents. Aldehydes are formed by the reaction of alkylmagnesium halides (Grignard reagents) with ethyl orthoformate or diethyl formamide [reactions (35) and (36)]. RMgX + HC(OCH2 CH3 )3 −→
C6H5CH2CI + Cl2
C6H5CHCl2 + HCI
H O−
3 RCH(OCH2 CH3 )2 −−−− −−→ RCHO (35)
(30)
Benzal chloride
O RMgX + HC N(CH3)2
Hydrolysis of benzal chloride gives benzaldehyde, as in reaction (31). This is one industrial process for C6H5CHCI2 + H2O
C6H5CHO + 2HCI
OMgX RCHN(CH3)2
(31)
Benzaldehyde
the production of benzaldehyde. Oxidation of a methyl group attached to a benzene ring. Certain substituted benzaldehydes, such as p-nitrobenzaldehyde, are readily prepared by oxidation of the corresponding toluene derivatives with chromic acid. Acetic anhydride is used to esterify the aldehyde groups and prevent further oxidation. Subsequent hydrolysis of the diacetate produces the aldehyde [reaction (32)].
H3O+
RCHO
(36)
Reimer-Tiemann synthesis. When phenols are treated with chloroform and strong alkali, o-hydroxybenzaldehydes (salicylaldehydes) are formed [reaction (37)].
OH
OH CHO + CHCI3
aq NaOH
R
(37) R
Phenols
o–Hydroxybenzaldehydes (salicylaldehydes)
O2NC6H4CH3 + O2 + 2(CH3CO)2O
p–Nitrotoluene
Furfural production. Cornstalks, corncobs, grain hulls, and similar farm wastes produce furfural upon heating with dilute sulfuric acid. This reaction (38), em-
Acetic anhydride O2NC6H4CH(OOCCH3)2
Pentosans
Diacetate
H2O
H2O
O2NC6H4CHO + 2HOOCH3
p–Nitrobenzaldehyde
Catalyst
Toluene
C6H5CHO + HCl Benzaldehyde
(33)
hydrogen in the presence of a specially prepared palladium catalyst to form aldehydes and hydrogen chloride, as in reaction (34). This reaction, known
Acid chloride
HCOH
heat acid
CH
O
Pentose HC
CH
HC
C
+ 3H2O
(38)
CHO
O Furfural
ployed on a large scale industrially, proceeds from pentosans found in the agricultural residues. Chloral production. Chloral is an important intermediate for the production of DDT. It is formed by chlorination of anhydrous ethanol [reaction (39)]. 2CH3CH2OH +
Pd
H2COH Acetic acid
Reduction of acid chlorides. Acid chlorides react with
RCOCl + H2
HCOH
(32)
One commercial synthesis of benzaldehyde involves the passing of a mixture of toluene and air over a heated metallic oxide catalyst at temperatures above 500◦C (930◦F) [reaction (33)]. C6H5CH3 + O2
HCOH
Ethanol
8Cl2 Chlorine 2Cl3CCHO + Chloral
1OHCl
(39)
Hydrogen chloride
RCHO + HCl
(34)
Aldehyde
as the Rosenmund synthesis, is a useful labora-
See CARBOHYDRATE; FORMALDEHYDE. Paul E. Fanta Bibliography. R. J. Fessenden and J. S. Fessenden, Organic Chemistry, 6th ed., 1998; T. W. G. Solomons, C. B. Fryhle, and T. G. Solomons, Organic Chemistry, 7th ed., 1999.
Aldosterone
Alder Any of the deciduous shrubs and trees of the genus Alnus in the birch family (Betulaceae). There are about 30 species, 10 in the United States, widespread in cool north temperate regions, and southward in the Andes of South America. They have a smooth gray bark, elliptical or ovate saw-toothed leaves in three rows, male flowers in long catkins mostly in early spring, and clusters of several dry hard ellipsoid blackish fruits 0.5–1 in. (1.2–2.5 cm) long. The fruits are conelike and present throughout the year. Alders are common in wet soils, such as stream borders. They often form thickets even beyond treelines and are pioneers on bare areas such as landslides, roadsides, and after fire or logging. As in legumes, the roots bear swellings or nodules containing nitrogenfixing bacteria which enrich the soil. See FAGALES; LEGUME; NITROGEN FIXATION.
Red alder (Alnus rubra) fruit and leaf.
Red alder (Alnus rubra; see illus.) is a small to large tree of the Pacific Northwest along the coast from southeast Alaska south to central California, and local in mountains of northern Idaho. It is the most common and most important hardwood in that region of conifers, occurring in pure stands and mixed with other species. This fast-growing, short-lived tree becomes 100–120 ft (30–36 m) tall on good sites. The wood is reddish or yellowish brown (whitish when freshly cut), moderately lightweight, moderately soft, and fine-textured. Principal uses are pulpwood, furniture, and cabinetwork. See FOREST AND FORESTRY; TREE. Elbert L. Little, Jr.
Aldosterone The principal mineralocorticoid (a steroid hormone that controls electrolyte and fluid balance) in humans. An adrenocortical steroid, aldosterone was first isolated in 1953 by S. A. Simpson and J. F. Tait in London. This hormone is synthesized from cholesterol in a process involving four different enzymes in the zona glomerulosa (outer zone) of the adrenal cortex. Aldosterone contains an equilibrium mixture of the aldehyde (see illustration) and the hemiacetal, the equilibrium favoring the latter. Normal serum concentrations range from 5 to 25 mg/dL. See ANDRENAL CORTEX; CHOLESTEROL; HORMONE; STEROID.
Biosynthesis. Direct stimulators of aldosterone biosynthesis include angiotensin II and III, ACTH, alpha-MSH, prolactin, vasopressin, potassium, hydrogen, ammonium, serotonin, histamine, and selected prostaglandins. The two major stimuli for aldosterone secretion are angiotensin II and potassium. Aldosterone synthesis is primarily regulated by the renin-angiotensin system. This system is activated by low sodium balance or decreased blood flow to the kidneys, where renin is produced. Renin is a proteolytic enzyme that converts the inactive angiotensinogen globulin to angiotensin I, which is then converted to angiotensin II by the angiotensin converting enzyme. See KIDNEY; THIRST AND SODIUM APPETITE. Physiologic activity. Aldosterone acts primarily at the renal distal convoluted tubule promoting the excretion of potassium and retention of sodium, thereby influencing extracellular volume homeostasis and blood pressure. See POTASSIUM; SODIUM. Alterations in Na+/K+ flux are seen between 30 and 90 minutes after aldosterone administration. This delayed response in activity is a nongenomic effect and represents binding to cytosolic steroid receptors, translocation to the nucleus, interaction with DNA, and finally genomic transcription and translation of effector proteins. Aldosterone stimulates the synthesis of epithelial sodium channels at the apical epithelial membrane of the epithelial cells in the kidney. Aldosterone increases the activity of the Na+/K+-ATPase. Mineralocorticoid receptors (the cytosolic steroid receptors) are present in epithelial cells in the kidney, bladder, gastrointestinal tract, sweat and salivary glands, smooth muscle, vascular endothelium, brain, and myocytes (muscle cells). Increased aldosterone secretion contributes to some forms of arterial hypertension, disorders of sodium retention and edema, and syndromes of potassium wasting (observed when intake of potassium is less than output in the urine). Chronic aldosterone excess in the presence of high salt intake causes cardiac fibrosis in experimental animals. Patients with chronic heart failure had a significant reduction in morbidity and mortality when given a mineralocorticoid receptor antagonist, spironolactone, in addition to traditional therapy. There are rapid (nongenomic) effects of aldosterone in smooth muscle, skeletal muscle, colonic
OH O HO
O
O Molecular structure of aldosterone, 11β,21-dihydroxy3,20-dioxo-pregn-4en-18al.
423
424
Alfalfa epithelial cells, and myocardial cells. These nongenomic effects have been linked to the development of increased systemic vascular resistance and therefore might participate in human hypertension and cardiovascular disease. See CARDIOVASCULAR SYSTEM; HYPERTENSION. Maria I. New; Alejandro Diaz Bibliography. J. W. Funder, Aldosterone action, Annu. Rev. Physiol., 55:115–130, 1993; D. N. Orth and W. J. Kovacs, The adrenal cortex, in J. D. Wilson et al. (eds.), William Textbook of Endocrinology, 9th ed., W. B. Saunders, Philadelphia, pp. 517–664, 1998; J. S. Willians and G. H. Willians, 50th aniversary of aldosterone, J. Clin. Endocrinol. Metab., 88:2364– 2372, 2003.
Alfalfa The world’s most valuable forage legume, Medicago sativa, also known also as lucerne, and less often as purple medic, medica, snail clover, median herb, Burgundy hay or clover, Chilean clover, and Burgoens Hoy. It is often referred to as the queen of forages. Alfalfa is produced worldwide on more than 32,000,000 hectares (80,000,000 acres). Seven countries—United States, Russia, Argentina, France, Italy, Canada, and China—account for 87% of the production area (see table). See LEGUME FORAGES; ROSALES.
Estimated hectarage in selected countries growing alfalfa Continent and country Europe Austria Bulgaria Czechoslovakia East Germany France Greece Hungary Italy Poland Romania Soviet Union (European) Spain Switzerland West Germany Yugoslavia North America Canada Mexico United States South America Argentina Chile Peru Asia Iran Soviet Union (Siberian) Turkey Africa Republic of South Africa Oceania Australia New Zealand *1
hectare = 2.5 acres.
Year
1982 1982 1983 1982 1983 1980 1982 1982 1981–1983 1981 1971 1981 1983 1983 1984
Hectarages*
12,630 399,000 200,000 190,000 566,000 198,700 337,500 1,300,000 258,000 400,000 3,375,000 332,600 6,000 31,000 337,000
1981 1982 1981
2,544,300 245,000 10,559,025
1981 1983 1981
7,500,000 60,000 120,000
1977 1971 1969
270,000 1,125,000 73,700
1985
300,000
1981–1982 1984
111,500 101,200
Description. Alfalfa is a herbaceous perennial legume. The seed, 0.04–0.08 in. (1–2 mm) long and wide, and 0.04 in. (1 mm) thick, is kidney-shaped and yellow to olive-green or brown. Plants produce a deep taproot that may have several branches. Two to twenty-five stems are borne from a single crown. Trifoliolate leaves are produced alternately on the stem. The flower is a simple raceme with several complete and perfect florets (Fig. 1), which may be various shades of purple, variegated, cream, yellow, and white. Pods are usually spirally coiled. Alfalfa is considered to be cross-pollinated. Under production conditions, pollination is by honeybees (Apis melifera), leafcutter bees (Megachile), and alkali bees (Nomia), depending on production location. Origin and domestication. Alfalfa originated in the Near East, in the area extending from Turkey to Iran and north into the Caucasus. Some scientists believe that alfalfa had two distinct centers of origin. The first center was in the mountainous region of Transcaucasia, and Asia Minor and the adjoining areas of northwestern Iran. The second area was in Central Asia. Domestication of alfalfa probably corresponded with the domestication of the horse. The oldest records trace to Turkey between 1400 and 1200 B.C., where alfalfa was fed to domestic animals during the winter and was regarded as a highly nutritious feed. The exact path of movement of alfalfa into the rest of the world is difficult to document. However, maritime trade was well developed in the eastern Mediterranean prior to 1400 B.C. and could have contributed to the spread of the crop. The Romans are known to have had alfalfa as early as the second century B.C. and are given credit for the development of advanced production practices. During this same period of time a Chinese expedition passing into Turkistan to obtain Iranian horses also obtained seed of alfalfa. With the fall of the Roman empire, alfalfa almost disappeared from Europe. During the eighteenth century alfalfa was taken from Europe to New Zealand, Australia, and the New World. There are two primary avenues of introduction of alfalfa into the United States. Spanish missionaries probably brought alfalfa from Chile and Mexico to Texas, New Mexico, Arizona, and California. The Chilean types were nonhardy, and thus climatic conditions prevented their movement into the northern United States and Canada. The second avenue of introduction was into Minnesota in 1857. A immigrant named Wendelin Grimm brought a few pounds of alfalfa seed (M. media = M. sativa × M. falcata) from Germany. Although his initial plantings were not very successful, by harvesting the seed from surviving plants he eventually developed a hardy type that became known as Grimm alfalfa. Establishment, management, and harvesting. Ideally, alfalfa should be planted in deep, well-drained, medium-textured soil with a pH of about 6.8. Because alfalfa is grown from a single planting for 3 or more years, land preparation is extremely important to successful establishment and production. In areas where irrigation is practiced, special attention must be paid to slope, levee construction, and drainage.
Alfalfa Soils should be able to supply the following kilograms of nutrients per megagram of hay produced (1 kg/Mg = 2 lb/ton): potassium, 40; calcium, 36; magnesium, 7; phosphorus, 6; sulfur, 5; and iron, chlorine, manganese, boron, zinc, copper, molybdenum, less than 0.5. In addition, the crop requires 123 lb (56 kg) of nitrogen, which is largely derived by nitrogen fixation, by Rhizobium meliloti. Seed should be inoculated with these bacteria, if there is any question of their presence in the soil. Planting rates vary between 16 and 45 kg/ha (14 and 40 lb/acre) depending on soil type and planting method. The higher rates are more common on heavy soils and when seeding is done by aircraft. Optimum soil temperatures for germination and seeding development are between 68 and 86◦F (20 and 30◦C). Dormant cultivars have lower optimum temperatures than nondormant cultivars. Forage should be harvested at one-tenth bloom or with 0.8– 1.2 in. (2–3 cm) of regrowth (Fig. 2). In northern production areas, one to three harvests are taken annually, whereas in areas such as the Imperial Valley of California up to ten harvests may be taken in a year. See NITROGEN CYCLE. Genetics and plant breeding. Alfalfa is an autotetraploid with a basic genomic number of 8 (2n = 4x = 32). However, diploid, hexaploid, and n = 7 species are represented in the genus. The frequency of quadrivalents is low, and so the frequency of double reduction is near zero. In the United States, prior to 1940 most of the improvement of alfalfa was by importation of germplasm adapted to production areas. Since the discovery of bacterial wilt disease (caused by Clavibacter michiganese ssp. insidiosum) considerable effort has been placed on the genetic improvement of the crop, with particular emphasis on the development of multiple pest and disease resistance. Phenotypic recurrent selection is the most common breeding method used for the improvement of alfalfa, but mass selection and various forms of genotypic recurrent selection are used. Cultivars released before 1960 were combinations of one to three germplasm sources, obtained from such areas as India, Chile, Peru, Africa, and Belgium, and primarily developed through public breeding programs. Subsequent to 1960, emphasis in public programs shifted to the development of breeding methods and germplasm. Most cultivars released since 1960 are populations (synthetics), proprietary, and combinations of several germplasm sources. A few commercial hybrids have been developed, and there is currently considerable interest in complementary strain crossing to produce new cultivars. See BREEDING (PLANT). Seed production and standards. Seed production is most efficient in hot and dry environments such as the San Joaquin Valley of California. Seed fields are usually planted on beds at 2–4 kg/ha (1.8–3.6 lb/ acre). Certified seed is produced under strict standards, which are monitored by state member organizations of the International Crop Improvement Association. These standards include the geographic area for seed production, the number of generations and
425
Fig. 1. Alfalfa stems at flowering stage. (Iowa State University Photo Service)
number of years of seed increase, isolation, freedom from volunteer plants and weeds, viability, and mechanical purity. Larry R. Teuber Diseases. Alfalfa is susceptible to at least 75 diseases caused by fungi, bacteria, mycoplasma-like
Fig. 2. Alfalfa of wide crown type at late bud or early bloom stage. (Iowa State University Photo Service)
426
Alfalfa
Fig. 3. Damaging effect of bacterial wilt of alfalfa (left), contrasted with healthy plant (right).
organisms, viruses, and nematodes. Most of these occur sporadically, and rarely in epidemic proportions; however, several are responsible for appreciable losses in forage yield and seed production. The best way to reduce disease losses in alfalfa is to grow locally adapted disease-resistant varieties. Root and crown diseases. Bacterial wilt is one of the best-known and most damaging diseases of alfalfa (Fig. 3). The bacteria enter the plants principally
Fig. 4. Fruiting structures of the Verticillium wilt fungus. Conidia (spores) are borne singly at the tips of conidiophores. (From J. H. Graham et al., A Compendium of Alfalfa Diseases, American Phytopathological Society, 1979)
through wounds in the roots and through cut ends of newly mowed stems. They are spread in the field by surface water and by tillage and harvesting machinery. Bacterial wilt causes stunting and reduced stand longevity. Resistant varieties have been developed for all alfalfa-growing regions of the United States. Verticillium wilt (Fig. 4), first identified in the United States in 1976, has been a damaging disease in Europe since the 1930s. It occurs in the northwestern United States and adjacent areas in Canada. Fusarium wilt, another fungus disease, occurs worldwide but is most severe in warm regions. Phytophthora root rot, caused by a water mold, occurs worldwide. It is favored by waterlogged soils resulting from excessive irrigation, abundant rainfall, or inadequate soil drainage. Resistant varieties are available. Other root and crown diseases include Sclerotinia crown rot, active during cool moist weather; Stagonospora root rot and leaf spot; and Fusarium and Rhizoctonia crown and root rots, more active during warm periods. Frequently, crown and root rots cannot be ascribed to a single pathogen, so the term disease complex is used to characterize the pathogens in combinations with various interacting agents (Fig. 5). These agents include stress factors such as winter damage, improper cultural and management practices, and root insects. No variety of alfalfa resistant to this complex of disorders has been developed so far. Stem and leaf diseases. Many microorganisms attack the stems and foliage of alfalfa. Some, such as the fungi that cause anthracnose and spring black stem, spread into the crown and upper part of the taproot, weakening or killing the plant. Anthracnose, favored by warm humid weather, is associated with the summer decline of alfalfa in the eastern United States. Resistance to anthracnose has been incorporated into some varieties; however, many of these varieties are susceptible to a second race of the anthracnose fungus which was discovered in 1978. Spring black stem occurs in most alfalfa-growing areas of the world and is particularly damaging in cooler humid regions. The fungus can attack any part of the plant, but the shiny black streaks on stems and irregularly shaped black spots on leaves are most conspicuous. Summer black stem is another disease of stems and leaves; its reddish to smoky-brown lesions appear when temperatures increase during summer. Several leaf diseases occur widely and cause losses through defoliation and a general weakening of the plant. Common leaf spot, yellow leaf blotch, Leptosphaerulina leaf spot, and Stemphylium leaf spot may develop so abundantly that they cause serious loss of leaves. Forage quality and yield are lowered. Some varieties have low levels of resistance. Early harvest may help reduce defoliation as well as spread of leaf and stem diseases. Some less prevalent but occasionally locally important foliar diseases are downy mildew, leaf rust, and bacterial leaf spot. Virus diseases. The alfalfa mosaic virus (AMV) complex consists of many strains that differ in virulence
Alfven ´ waves with those caused by infectious (biotic) agents. In addition, abiotic agents often influence the severity of diseases caused by biotic agents. See PLANT PATHOLOGY. J. H. Graham Bibliography. D. K. Barnes et al., Alfalfa germplasm in the United States: Genetic vulnerability, use improvement and maintenance, USDA Tech. Bull. 1571, 1977; C. H. Hanson (ed.), Alfalfa Science and Technology, American Society of Agronomy, 1972; A. A. Hanson, R. R. Hill, and D. K. Barnes (eds.), Alfalfa and Alfalfa Improvement, American Society of Agronomy, 1988; D. Smith, Forage Management in the North, 1975.
Alfven ´ waves
Fig. 5. Crown rot of alfalfa. It may be caused by several fungi in combination with stress factors such as winter injury. (From J. H. Graham et al., A Compendium of Alfalfa Diseases, American Phytopathological Society, 1979)
and symptoms. The classic symptom is yellow mottling of leaves. Other symptoms are stunting, contortions of leaves and petioles, and, in some instances, death of plants. However, most AMV-infected plants in an alfalfa stand never show symptoms. Exact information on disease losses is lacking. Because the virus is seedborne, seed produced in areas of low virus incidence should be used. Other viruses of alfalfa include alfalfa enation, transient streak, and alfalfa latent viruses. See PLANT VIRUSES AND VIROIDS. Mycoplasmalike diseases. Witches’-broom, long believed to be a leafhopper-transmitted virus, is now thought to be caused by a mycoplasmalike microorganism. The disease is of minor importance in the United States but is widespread in parts of semiarid Australia. Nematodal diseases. Plant-parasitic nematodes of alfalfa not only damage plants directly but also aid in the transmission and entry of other pathogens. The alfalfa stem nematode is most damaging, especially in the western United States. Infected crown buds and stems are swollen and stunted and may be killed. Reuse of irrigation water is a major pathway for dissemination. Some resistant varieties have been developed. The root-knot nematodes (Meloidogyne spp.) and the root-lesion nematodes (Pratylenchus spp.) are widespread and attack many crop plants and weeds. There is usually less damage to the alfalfa crop than to more susceptible crops that follow in a rotation. Noninfectious diseases. Alfalfa frequently shows disease symptoms not associated with pathogenic or parasitic organisms. These noninfectious or abiotic diseases are caused by the lack or excess of plant nutrients, extreme temperatures, air pollutants, pesticides, harvesting machinery, and other factors. Symptoms caused by these agents can be confused
Propagating oscillations in electrically conducting fluids or gases in which a magnetic field is present. Magnetohydrodynamics deals with the effects of magnetic fields on fluids and gases which are efficient conductors of electricity. Molten metals are generally good conductors of electricity, and they exhibit magnetohydrodynamic phenomena. Gases can be efficient conductors of electricity if they become ionized. Ionization can occur at high temperatures or through the ionizing effects of high-energy (usually ultraviolet) photons. A gas which consists of free electrons and ions is called a plasma. Most gases in space are plasmas, and magnetohydrodynamic phenomena are expected to play a fundamental role in the behavior of matter in the cosmos. See MAGNETISM; PLASMA (PHYSICS). Waves are a particularly important aspect of magnetohydrodynamics. They transport energy and momentum from place to place and may, therefore, play essential roles in the heating and acceleration of cosmical and laboratory plasmas. A wave is a propagating oscillation. If waves are present, a given parcel of the fluid undergoes oscillations about an equilibrium position. The parcel oscillates because there are restoring forces which tend to return it to its equilibrium position. In an ordinary gas, the only restoring force comes from the thermal pressure of the gas. This leads to one wave mode: the sound wave. If a magnetic field is present, there are two additional restoring forces: the tension associated with magnetic field lines, and the pressure associated with the energy density of the magnetic field. These two restoring forces lead to two additional wave modes. Thus there are three magnetohydrodynamic wave modes. However, each restoring force does not necessarily have a unique wave mode associated with it. Put another way, each wave mode can involve more than one restoring force. Thus the usual sound wave, which involves only the thermal pressure, does not appear as a mode in magnetohydrodynamics. See SOUND; WAVE MOTION IN FLUIDS. The three modes have different propagation speeds, and are named fast mode (F), slow mode (S), and intermediate mode (I). The intermediate mode is sometimes called the Alfv´en wave, but some
427
428
Alfven ´ waves scientists refer to all three magnetohydrodynamic modes as Alfv´en waves. The intermediate mode is also called the shear wave. Some scientists give the name magnetosonic mode to the fast mode. Basic equations. The magnetohydrodynamic wave modes are analyzed by using the magnetohydrodynamic equations for the motion of a conducting fluid in a magnetic field, combined with Maxwell’s equations and Ohm’s law. See MAXWELL’S EQUATIONS. In studies of waves, the most important of the magnetohydrodynamic equations is called the momentum equation. It expresses how fluid parcels accelerate in response to the various forces acting on the fluid. The essential forces come from the thermal pressure in the fluid, and from the tension and pressure forces associated with the magnetic field. The magnetic field can exert forces on the fluid only when an electric current flows in the fluid; thus magnetohydrodynamic waves can exist only in an electrically conducting fluid. See FLUID FLOW. Two other equations express the conservation of mass of the fluid and the conservation of entropy. However, entropy is conserved only when viscosity, heat conduction, and electrical resistivity are ignored. See CONSERVATION OF MASS; ENTROPY. The fourth equation is Ohm’s law, which expresses the relationship between the electric current and the electric field in the fluid. If the fluid has no electrical resistivity, any nonzero electric field would lead to an infinite current, which is physically unacceptable. Thus the electric field “felt” by any fluid parcel must be zero. However, detailed analysis of the motions of electrons and ions reveals that this statement is valid only if the wave frequency is much less than the cyclotron frequency of the plasma ions; thus magnetohydrodynamic waves are very low-frequency waves. See OHM’S LAW; PARTICLE ACCELERATOR; RELATIVISTIC ELECTRODYNAMICS. It is possible to combine Ohm’s law with Faraday’s law of induction. The resultant equation is called the magnetohydrodynamic induction equation, which is the mathematical statement of the “frozen-in” theorem. This theorem states that magnetic field lines can be thought of as being frozen into the fluid, with the proviso that the fluid is always allowed to slip freely along the field lines. The coupling between the fluid and the magnetic field is somewhat like the motion of a phonograph needle in the grooves of a record: the needle is free to slip along the grooves (the field lines) while otherwise being constrained from crossing the grooves. It is this coupling between the fluid and the magnetic field which makes magnetohydrodynamic waves possible. The oscillating magnetic field lines cause oscillations of the fluid parcels, while the fluid provides a mass loading on the magnetic field lines. This mass loading has the effect of slowing down the waves, so that they propagate at speeds much less than the speed of light (which is the propagation speed of waves in a vacuum). See ELECTROMAGNETIC RADIATION; FARADAY’S LAW OF INDUCTION; LIGHT. The final equation is Amp`ere’s law, which states that an electric current always produces a magnetic field. See AMPE`RE’S LAW; BIOT-SAVART LAW.
Thus there are six equations for six unknowns (the pressure, density, and velocity of the fluid; the electric current; and the electric and magnetic fields). Six equations are in principle sufficient to determine six unknowns. Linearization of equations. Unfortunately, the six equations are too difficult to be of much use because some of them are nonlinear; that is, they contain products of the quantities for which a solution is sought. Nonlinear magnetohydrodynamics is still only in its infancy, and only a few specialized solutions are known. In order to get solvable equations, scientists accept the limitation of dealing with smallamplitude waves and linearize the equations, so that products of the unknowns are removed. Fortunately, much can still be learned from this procedure. The linearization consists of four steps: (1) All quantities are split into two parts: a background part (denoted by the subscript 0) which is constant in space and time, and a wave part (denoted by the prefix δ) which oscillates in space and time; for example, B = B0 + δB and v = v0 + δv, where B denotes the vector magnetic field and v denotes fluid flow velocity. (2) These definitions are inserted into the six fundamental equations. (3) It is assumed that the waves are of such small amplitude that products of the wave parts can be ignored. (4) It is assumed that there is no background flow. (This last step simplifies the mathematics without loss of generality.) The resulting equations have solutions which are harmonic in time and space, that is, which vary as exp(ik · r − iωt). Here r, is spatial position, ω is (angular) frequency, and k is called the wave vector. The magnitude of the wave vector represents how rapidly the wave oscillates in space, and its direction represents the direction of propagation. The magnitude of the angular frequency represents how rapidly the wave oscillates in time. Since the equations are linear, more general solutions can be constructed by adding together an arbitrary number of harmonic solutions. See HARMONIC MOTION. Dispersion relation. The linearized equations have solutions only when ω and k are related to one another in a special way. The equation which expresses this relationship between ω and k is called the dispersion relation. It is a fundamental relationship in the theory of waves. See DISPERSION RELATIONS. The dispersion relation for magnetohydrodynamic waves is given by Eq. (1). The quantity M02 is given by Eq. (2). The quantity k is the component of k (ω2 − k2 v2A )(k2⊥ + M02 ) = 0 M02 =
(k2 cs2 − ω2 )(k2 v2A − ω2 ) (cs2 + v2A )(k2 cT2 − ω2 )
(1) (2)
parallel to B0 while k⊥ is the component perpendicular to B0. Equation (2) contains several velocities which are fundamental velocities in magnetohydrodynamics. The quantity υ A is called the Alfv´en speed. The Alfv´en speed is equivalent to the speed at which a wave travels along a taut string, except that the tension of the string is replaced by the tension in
Alfven ´ waves the magnetic field lines. The quantity cs is the usual speed of sound in a gas, and cT is called the tube speed, given by Eq. (3).
1.5
F
(4)
0.5
0
0°
I
I
S
S
30°
60°
90°
0°
90°
(b)
Fig. 1. Phase speed as a function of propagation direction for the fast (F), slow (S), and ´ speed intermediate (I) modes. The ordinate is phase velocity Vph, normalized to the Alfven VA. The abscissa is the angle θ between the wave propagation of the direction of the magnetic field B0. (a) Cs/VA = 0.75. (b) Cs/VA = 1.5.
lated (Fig. 2) and in fact satisfy Eq. (7). (The symδv/vA = −sgn(k · B0 )δB/B0
(7)
bol sgn means the algebraic sign of the quantity in parentheses.) Equation (7) has been used to identify the presence of the intermediate mode in the solar wind. There, spacecraft-borne instruments have measured the magnetic field and plasma velocity directly. The fluctuations have been found to obey Eq. (7) fairly closely. It appears that intermediate waves are rather copious in the solar wind. Further investigations have revealed that the waves usually propagate away from the Sun. The Sun may thus be an emitter of these waves. However, the waves that are found in the solar wind also exhibit many properties of turbulence; thus, nonlinear effects, such as interactions between different waves, are important. See SOLAR WIND; SUN. i −1δj
+ δB
k
θ B0 δV
(6)
quantity k⊥ does not appear because the motions on neighboring field lines do not communicate with one another. This is a consequence of the fact that the motions are shears, so that neighboring field lines never bump together. The fluctuating parts of the fluid flow velocity and the magnetic field, δv and δB, are closely re-
60°
30° θ
(a)
(5)
change character depending on whether cs/vA is less than or greater than 1. See PHASE VELOCITY. All three modes can propagate along the magnetic field, while only the fast mode can propagate across the magnetic field. For propagation along the field (θ = 0), two of the modes have vph = vA, while the third has vph = cs; the latter mode is essentially a sound wave unaffected by the magnetic field, while the other two modes are intrinsically magnetic. The intermediate mode has vph = vA cos θ ; this mode is always magnetic in character, and thus the sound speed does not affect its phase speed. In all other cases the fast and slow modes have a mixed character, and all three restoring forces combine in a complicated way. The essential properties of the modes are summarized below. Intermediate mode. The velocity and magnetic field fluctuations in this mode are perpendicular to the plane containing k and B0. This means that the motions are pure shears. There is no compression of the plasma. This is why the sound speed does not appear in the dispersion relation. The tension in the magnetic field lines is the only restoring force involved in the propagation of the wave. This mode is therefore closely analogous to the propagation of waves on a string. The dispersion relation is given by Eq. (6). The ω2 = k2 v2A
1.0
θ
This quantity can be obtained by solving Eq. (1). The three modes are designated fast, intermediate, and slow, according to their phase speeds (Fig. 1). The phase speed depends on the angle θ between the wave propagation direction and the direction of B0. Thus θ satisfies Eq. (5). The graphs of phase speed tan θ = k⊥ /k
F
(3)
Equation (1) is a cubic equation for ω2. The three roots represent the three magnetohydrodynamic wave modes. For each of the three values of ω2, there are two values of ω (one positive and one negative) which correspond to waves propagating in opposite directions. A particularly useful quantity is the propagation speed, or phase speed, of the wave, given by Eq. (4). vph = ω/k
2.0
V ph / V A
cT2 = cs2 v2A /(cs2 + v2A )
429
δE Fig. 2. Vector relationship between the various fluctuating wave quantities for the intermediate mode. Symbols next to δB and δv indicate that these quantities are pointing into and out of the paper, respectively. √ The electric current δj must be multiplied by i−1(or 1/ − 1) because the current is ◦ 90 out of phase with the other quantities.
430
Alfven ´ waves Waves propagate energy. Since the intermediate mode does not involve the thermal properties of the plasma, its energy flux density is given by the time average of the electromagnetic Poynting flux. For the intermediate mode, the wave energy flux density is always parallel or antiparallel to B0. Thus the intermediate mode always channels energy along the magnetic field. In fact, the energy flux density of the wave consists of the wave’s kinetic energy density plus the wave’s magnetic energy density, propagating along the magnetic field at speed vA. See POYNTING’S VECTOR. Because these waves channel energy along magnetic fields, they may be responsible for the observed fact that cosmical plasmas are strongly heated in the presence of magnetic fields. For example, in the solar wind the wave energy is observed to decrease at greater distances from the Sun, and the amount of wave energy that is lost is about the same as the heat energy that is gained by the solar-wind plasma. Thus, the solar wind is an example of a cosmical plasma that is heated by waves. Waves may also accelerate the solar wind through an effect which is closely analogous to the radiation pressure exerted by light waves. They may also heat the solar atmosphere. Results from the Solar and Heliospheric Observatory (SOHO) indicate that protons and heavier ions are strongly heated close to the Sun, in the region where the solar wind is formed. The heating is strongest in the directions perpendicular to B0. The high-frequency intermediate mode (also called the ion-cyclotron wave) is capable of producing this kind of heating. However, it has not yet been proven that the waves are present at the Sun. It is even less certain whether waves heat the atmospheres of other stars, or other astrophysical objects. Intermediate waves have also been proposed as a means of heating plasmas in controlled thermonuclear fusion devices. See NUCLEAR FUSION. Fast mode. The fast mode is difficult to analyze. However, in many cosmical and laboratory plasmas, νA2 is much greater than CS2 . This is the strongmagnetic-field case. In that case the fast mode is more easily understood. From Eq. (1), the fast-mode dispersion relation is approximately given by Eq. (8). The direction of k ω2 /k2 ≈ v2A
is much greater than CS2 . The slow mode in a strong field is equivalent to sound waves which are guided along the strong magnetic field lines. The strong magnetic field lines can be thought of as a set of rigid pipes which allow free fluid motion along the pipes, but which restrict motion in the other two directions. The motions on the individual pipes are not coupled together, and thus the slow mode is analogous to the sound waves on a set of independent organ pipes. The slow mode channels energy along the magnetic field. Because the sound speed is small, by assumption, the slow mode transmits energy less effectively than the fast or intermediate modes. Nonlinear effects. Only small-amplitude waves have been considered. Real waves have finite amplitude, and nonlinear effects can sometimes be important. One such effect is the tendency of waves to steepen, ultimately forming magnetohydrodynamic shock waves and magnetohydrodynamic discontinuities. There is an abundance of magnetohydrodynamic discontinuities in the solar wind. Nonlinear studies have shown these are the result of the nonlinear steepening of the intermediate waves. Nonlinear studies have also shown that the intermediate mode will naturally reach a state where the magnetic field can have large fluctuations of its direction, but maintain a constant field strength. Thus the tip of the magnetic field vector moves on the surface of a sphere. This state of spherical polarization is in fact commonly observed in the solar wind. See NONLINEAR ACOUSTICS; SHOCK WAVE. It is also possible that waves can degenerate into turbulence. There are indications that this too happens in the solar wind. The observed power spectra show that fluctuations are present with a very broad range of frequencies and wave-vector magnitudes, which is a property of turbulence. At greater distances from the Sun, the outward-propagating waves are less dominant; they are apparently getting scrambled because of turbulent processes. See TURBULENT FLOW.
δv k
(8)
does not appear in Eq. (8). Thus the waves propagate isotropically. Fast waves are compressive, and the magnetic field strength fluctuates as well (Fig. 3). Thus fast waves are governed by the two restoring forces associated with the tension and pressure in the magnetic field. Because the magnetic field is strong, the wave energy flux density is mainly due to the time average of the Poynting flux. For this mode, it is evident that the energy flux is along k. Thus the fast mode can propagate energy across the magnetic field. Slow mode. Like the fast mode, the slow mode is difficult to study in general, and the discussion will again be confined to strong magnetic fields, so that v2A
θ +
i − 1 δj
B0 δE
δB Fig. 3. Vector relationships between the various fluctuating wave quantities for the fast mode in a strong magnetic field, where vA2 is much greater than cA2.
Algae Dissipation. Viscosity, heat conduction, and electrical resistivity have all been ignored in the above discussion. These effects will all tend to dissipate the wave energy and heat the plasma. See CONDUCTION (HEAT); ELECTRICAL RESISTIVITY; SOUND ABSORPTION; VISCOSITY. In some plasmas, the electrons and ions almost never collide with one another. In such nearly collisionless plasmas, the fast and slow modes can dissipate by the process of Landau damping. Shock formation and turbulence also lead to wave dissipation. Surface waves. Only waves in a spatially uniform background have been considered. While the analysis of magnetohydrodynamic waves in a nonuniform background is complicated, it is possible to consider an extreme limit, in which the background is uniform except at certain surfaces where it changes discontinuously. Surfaces can support magnetohydrodynamic waves, which are in some respects similar to waves on the surface of a lake. These waves may play important roles in heating cosmical and laboratory plasmas. See MAGNETOHYDRODYNAMICS; WAVE MOTION IN LIQUIDS. Joseph V. Hollweg Bibliography. R. L. Carovillano and J. M. Forbes (eds.), Solar-Terrestrial Physics, 1983; F. F. Chen, Introduction to Plasma Physics, 1984; J. V. Hollweg, Computer Phys. Rep., 12:205–232, 1990; E. R. Priest, Solar Magnetohydrodynamics, 1982; R. Schwenn and E. Marsch (eds.), Physics of the Inner Heliosphere, vol. 2, 1991.
Algae An informal assemblage of predominantly aquatic organisms that carry out oxygen-evolving photosynthesis but lack specialized water-conducting and foodconducting tissues. They may be either prokaryotic (lacking an organized nucleus) and therefore members of the kingdom Monera, or eukaryotic (with an organized nucleus) and therefore members of the kingdom Plantae, constituting with fungi the subkingdom Thallobionta. They differ from the next most advanced group of plants, Bryophyta, by their lack of multicellular sex organs sheathed with sterile cells and by their failure to retain an embryo within the female organ. Many colorless organisms are referable to the algae on the basis of their similarity to photosynthetic forms with respect to structure, life history, cell wall composition, and storage products. The study of algae is called algology (from the Latin alga, meaning sea wrack) or phycology (from the Greek phykos, seaweed). See BRYOPHYTA; PLANT KINGDOM; THALLOBIONTA. General form and structure. Algae range from unicells 1–2 micrometers in diameter to huge thalli [for example, kelps often 100 ft (30 m) long] with functionally and structurally distinctive tissues and organs. Unicells may be solitary or colonial, attached or free-living, with or without a protective cover, and motile or nonmotile. Colonies may be irregular or with a distinctive pattern, the latter type being flagellate or nonmotile. Multicellular algae form pack-
ets, branched or unbranched filaments, sheets one or two cells thick, or complex thalli, some with organs resembling roots, stems, and leaves (as in the brown algal orders Fucales and Laminariales). Coenocytic algae, in which the protoplast is not divided into cells, range from microscopic spheres to thalli 33 ft (10 m) long with a complex structure of intertwined siphons (as in the green algal order Bryopsidales). Classification. Sixteen major phyletic lines (classes) are distinguished on the basis of differences in pigmentation, storage products, cell wall composition, flagellation of motile cells, and structure of such organelles as the nucleus, chloroplast, pyrenoid, and eyespot. These classes are interrelated to varying degrees, the interrelationships being expressed by the arrangement of classes into divisions (the next-higher category). Among phycologists there is far greater agreement on the number of major phyletic lines than on their arrangement into divisions. Superkingdom Prokaryotae Kingdom Monera Division Cyanophycota (= Cyanophyta, Cyanochloronta) Class Cyanophyceae, blue-green algae Division Prochlorophycota (= Prochlorophyta) Class Prochlorophyceae Superkingdom Eukaryotae Kingdom Plantae Subkingdom Thallobionta Division Rhodophycota (= Rhodophyta, Rhodophycophyta) Class Rhodophyceae, red algae Division Chromophycota (= Chromophyta) Class: Chrysophyceae, golden or goldenbrown algae Prymnesiophyceae (= Haptophyceae) Xanthophyceae (= Tribophyceae), yellow-green algae Eustigmatophyceae Bacillariophyceae, diatoms Dinophyceae, dinoflagellates Phaeophyceae, brown algae Raphidophyceae, chloromonads Cryptophyceae, cryptomonads Division Euglenophycota (= Euglenophyta, Euglenophycophyta) Class Euglenophyceae Division Chlorophycota (= Chlorophyta, Chlorophycophyta) Class: Chlorophyceae, green algae Charophyceae, charophytes Prasinophyceae Placing more taxonomic importance on motility than on photosynthesis, zoologists traditionally have considered flagellate unicellular and colonial algae as protozoa, assigning each phyletic line the rank of order. See BACILLARIOPHYCEAE; CHAROPHYCEAE; CHLOROPHYCEAE; CHRYSOPHYCEAE; CRYPTOPHY-
431
432
Algae CEAE; CYANOPHYCEAE; DINOPHYCEAE; EUGLENOPHYCEAE; EUKARYOTAE; EUSTIGMATOPHYCEAE; PHAEOPHYCEAE; PRASINOPHYCEAE; PROCHLOROPHYCEAE; PROKARYOTAE; PROTOZOA; PRYMNESIOPHYCEAE; RAPHIDOPHYCEAE; RHODOPHYCEAE; THALLOBIONTA; XANTHOPHYCEAE.
Cell covering. Although some unicellular algae are naked or sheathed by mucilage or scales, most are invested with a covering (wall, pellicle, or lorica) of
diverse composition and construction. These coverings consist of at least one layer of polysaccharide (cellulose, alginate, agar, carrageenan, mannan, or xylan), protein, or peptidoglycan that may be impregnated or encrusted with calcium carbonate, iron, manganese, or silica. They are often perforated and externally ornamented (Fig. 1a, d, i, j, and k). Diatoms have a complex wall composed almost entirely of silica. In multicellular and coenocytic algae,
Fig. 1. Representative fresh-water algae. (a) Phacotus (Chlorophyceae), motile unicell with a lorica; (b) Golenkinia (Chlorophyceae), nonmotile planktonic unicell; (c) Vaucheria (Xanthophyceae), siphonous coenocyte; (d) Phacus (Euglenophyceae), uniflagellate unicell; (e) Tetraedron (Chlorophyceae), nonmotile planktonic unicell; (f) Botrydium (Xanthophyceae), terrestrial spheroidal coenocyte; (g) Microcystis (Cyanophyceae), coccoid colony; (h) Actidesmium (Chlorophyceae), irregular colony; (i) Tetragoniella (Xanthophyceae), nonmotile unicell; (j) Euastrum (Chlorophyceae), desmid showing semicells; (k) Navicula (Bacillariophyceae), unicell with siliceous wall; (l) Crucigenia (Chlorophyceae), regular nonmotile colony; (m) Cystodinium (Dinophyceae), nonmotile cystlike dinoflagellate: (n) Pandorina (Chlorophyceae), regular motile colony; (o) Asterococcus (Chlorophyceae), unicell in mucilaginous sheath; (p) Ceratium (Dinophyceae), motile dinoflagellate; (q) Stigeoclonium (Chlorophyceae), branched filament; (r) Chroococcus (Cyanophyceae), unicell, solitary or in irregular colony; (s) Pediastrum (Chlorophyceae), regular nonmotile colony; (t) Characium (Chlorophyceae), epiphytic unicell; (u) Gonyostomum (Raphidophyceae), biflagellate unicell; (v) Merismopedia (Cyanophyceae), platelike colony; (w) Wollea (Cyanophyceae), one trichome from a gelatinous colony; (x) Ulothrix (Chlorophyceae), unbranched filament; (y) Oscillatoria (Cyanophyceae), trichome lacking sheath.
Algae most reproductive cells are naked, but vegetative cells have walls whose composition varies from class to class. See CELL WALLS (PLANT). Cytoplasmic structure. Prokaryotic algae lack membrane-bounded organelles. Eukaryotic algae have an intracellular architecture comparable to that of higher plants but more varied. Among cell structures unique to algae are contractile vacuoles in some freshwater unicells, gas vacuoles in some planktonic blue-green algae, ejectile organelles in dinoflagellates and cryptophytes, and eyespots in motile unicells and reproductive cells of many classes. Nuclei of different taxonomic groups vary in the degree of nuclear breakdown during division, the presence of histone around the deoxyribonucleic acid (DNA), details of spindle formation and elongation, and chromosome behavior. Chromosome numbers vary from n = 2 in some red and green algae to n ≥ 300 in some dinoflagellates. The dinoflagellate nucleus is in some respects intermediate between the chromatin region of prokaryotes and the nucleus of eukaryotes and is termed mesokaryotic. Some algal cells characteristically are multinucleate, while others are uninucleate. Chloroplasts, which always originate by division of preexisting chloroplasts, have the form of plates, ribbons, disks, networks, spirals, or stars and may be positioned centrally or along the cell wall. Photosynthetic membranes (thylakoids) are arranged in distinctive patterns and contain pigments diagnostic of individual classes. Chloroplast DNA may be dispersed throughout the organelle or (in Chromophycota) distributed in a ring. Pyrenoids— proteinaceous bodies around which carbohydrate is stored—occur in the chloroplasts of at least some species of all classes. They have been shown to contain the same carboxylating enzyme that is present in higher plants with C3 metabolism. See CELL (BIOLOGY); CELL PLASTIDS; CHROMOSOME; DEOXYRIBONUCLEIC ACID (DNA); PHOTOSYNTHESIS; PLANT CELL. Movement. In all classes of algae except Prochlorophyceae, there are cells that are capable of movement. The slow, gliding movement of certain bluegreen algae, diatoms, and reproductive cells of red algae presumably results from extracellular secretion of mucilage. Ameboid movement, involving pseudopodia, is found in certain Chrysophyceae and Xanthophyceae. An undulatory or peristaltic movement occurs in some Euglenophyceae. The fastest movement is produced by flagella, which are borne by unicellular algae and reproductive cells of multicellular algae representing all class– es except Cyanophyceae, Prochlorophyceae, and Rhodophyceae. Most motile cells have two flagella, inserted apically or laterally, but some have three, four, eight, or a crown of flagella. Those that appear to be uniflagellate probably have a second, nonemergent flagellum. The basic structural feature of a flagellum—two central microtubules surrounded by nine doublets—is constant throughout the algae (as in all other plants and animals) except in male gametes of diatoms, in which the central pair is absent. Flagella may be decorated
with hairs, spines, or scales, the details being of diagnostic value at various taxonomic levels. Details of the flagellar basal region are extremely variable and are considered to be of phylogenetic significance. Biflagellate cells may be isokont (with identical flagella) or heterokont (with morphologically different flagella). Flagellate cells with eyespots are phototactic, moving into areas of optimal light intensity. Internal movement also occurs in algae in the form of cytoplasmic streaming and light-induced orientation of chloroplasts. See CELL MOTILITY; CILIA AND FLAGELLA. Cell division and growth. Cell division is the process by which unicellular algae multiply and multicellular algae increase in size and complexity. Even within a single class, such as the green algae, details of cell division vary with respect to involvement and interaction of nuclei, microtubules, vesicles, and plasmalemma in wall formation. Random cell division results in diffuse growth, as in a filament or sheet. Thalli with differentiated organs, however, require the localization of cell division in discrete areas or tissues (meristems), which may be apical, basal, marginal, or intercalary. Growth in coenocytic algae entails expansion or elongation of the wall and repeated division of the nucleus and other organelles, but without formation of internal walls. See CELL DIVISION. Reproduction and life histories. Sexual reproduction is unknown in prokaryotic algae and in three classes of eukaryotic unicells (Eustigmatophyceae, Cryptophyceae, and Euglenophyceae), in which the production of new individuals is by binary fission. In sexual reproduction, which is found in all remaining classes, the members of a copulating pair of gametes may be morphologically indistinguishable (isogamous), morphologically distinguishable but with both gametes motile (anisogamous), or differentiated into a motile sperm and a relatively large nonmotile egg (oogamous). Gametes may be formed in undifferentiated cells or in special organs (gametangia), male (antheridia) and female (oogonia). (The highly modified gametangia of red algae are called spermatangia and carpogonia, respectively.) When gametangia are multicellular, every cell becomes fertile, sterile protective cells being absent. Sex hormones (attractants and determinants) have been demonstrated in a few brown and green algae. Sexual reproduction may be replaced or supplemented by asexual reproduction, in which special cells (spores) capable of developing directly into a new alga are formed in undifferentiated cells or in distinctive organs (sporangia). Spores may be motile or nonmotile, and may germinate immediately upon release or may secrete a thick wall to protect themselves during a period unfavorable to germination. Many multicellular algae reproduce vegetatively by fragmentation or by the formation of propagules. See REPRODUCTION (PLANT). Life histories range from binary fission in some unicells to an elaborate sequence of four multicellular somatic phases in many red algae, involving
433
434
Algae highly specialized reproductive systems. The different life history patterns are defined in terms of the position of syngamy and meiosis in relation to somatic phases. The three basic types of life history described below may be characterized as having zygotic meiosis, sporic meiosis, and gametic meiosis, respectively. In exceptional cases, meiosis regularly occurs in a somatic tissue, resulting in chimeralike thalli and a complicated life history. In fresh-water algae, sexual reproduction is generally correlated with the onset of adverse environmental conditions (drying up of the habitat in summer or freezing in winter). The zygote resulting from gametic union secretes a thick, protective wall, transforming itself into a resting stage (zygospore). At the return of favorable growing conditions, the zygote germinates by undergoing meiosis and giving rise usually to four haploid spores, each of which produces a new thallus. The habitat of marine algae, by contrast, undergoes seasonal changes more gradually, and the zygote usually develops immediately into a diploid somatic phase. Adverse conditions of light and temperature are endured, not by thick-walled spores, but by filamentous or crustose stages that receive protection from predation and physical stress by nestling in crevices. In green and brown marine algae, the diploid somatic phase is autonomous, but in red algae, which have two successive diploid phases, only the second is autonomous, the first remaining an integral part of the parent gamete-producing plant (gametophyte). The autonomous diploid plant (sporophyte) in turn produces spores by meiosis. The haploid spores (meiospores) then give rise directly to unisexual (dioecious) or bisexual (monoecious) gametophytes, which produce gametes to complete the cycle. This life history, which involves an alternation or sequence of somatic phases correlated with changes in ploidy level, is often inappropriately called alternation of generations. The haploid and diploid somatic phases may be similar to one another (isomorphic) or dissimilar in varying degrees (heteromorphic). Extreme cases are presented by kelps, in which the sporophyte often reaches a length of 100 ft (30 m) while the gametophytes are microscopic. A few algae (Fucales, some Bryopsidales) have a life history in which there is only one somatic phase (a bisexual plant) or two expressions (male and female) of one somatic phase, which are diploid. The formation of gametes entails meiosis, and the zygote develops directly into a new diploid thallus. Culture studies supplemented by field studies have revealed that some algae are extremely plastic with regard to life history, being able to bypass one or another phase by recycling (repeating), apomeiosis (short-circuiting of sexual reproduction), pedogenesis, parthenogenesis, and vegetative multiplication, and to produce special microscopic (usually filamentous) stages in response to seasonal, long-term, or sporadic changes in the environment. Some species of benthic marine algae have different life histories in different parts of their geographic
range, determined probably by temperature and light regimes. Metabolism. Most algae are autotrophic, obtaining energy and carbon through photosynthesis. All photosynthetic algae liberate oxygen and use chlorophyll a as the primary photosynthetic pigment. Secondary (accessory) photosynthetic pigments, which capture light energy and transfer it to chlorophyll a, include chlorophyll b (Prochlorophyceae, Euglenophyceae, Chlorophycota), chlorophyll c (Chromophycota), fucoxanthin among other xanthophylls (Chromophycota), and phycobiliproteins (Cyanophyceae, Rhodophyceae, Cryptophyceae). Other carotenoids, especially β-carotene, protect the photosynthetic pigments from oxidative bleaching. Except for different complements of accessory pigments (resulting in different action spectra), photosynthesis in algae is identical to that in higher plants. Carbon is predominantly fixed through the C3 pathway. See CAROTENOID; CHLOROPHYLL. Obligately or facultatively heterotrophic algae occur in most classes. These algae are parasites (on other algae, higher plants, and a broad spectrum of animals including humans) or free-living osmotrophs or phagotrophs. The chemical products of photosynthesis in most algae are stored as one of two polymers of glucose— starch similar to that in higher plants, in which the glucose units have α-1,4-linkages, or laminaran, chrysolaminaran, and paramylon, with β-1,3linkages. Other storage products (sugars, alcohols, and lipids) are less common, but may be abundant in certain taxonomic groups. Phosphorus in the form of polyphosphate and nitrogen in the form of proteins or peptides are stored by some algae. The source of carbon for most photosynthetic algae is carbon dioxide (CO2), but some can use bicarbonate. Many photosynthetic algae are also able to use organic substances (such as hexose sugars and fatty acids) and thus can grow in the dark or in the absence of CO2. Colorless algae obtain both energy and carbon from a wide variety of organic compounds in a process called oxidative assimilation. Most algae obtain nitrogen for protein synthesis from nitrate or ammonium, but some unicells can use organic sources such as urea or amino acids. Some blue-green algae—with few exceptions only those with heterocysts—can meet their requirements by fixing atmospheric nitrogen. See NITROGEN FIXATION. Algae, like higher plants, require a certain complement of elements for successful growth and reproduction. In addition to the elements required by higher plants, at least some algae require vanadium, silicon, iodine, or sodium. Unlike higher plants, many algae need an external source of certain vitamins (vitamin B12, thiamin, and biotin, singly or in combination). These vitamins are usually available in natural waters as a result of the metabolic activity of bacteria, but they may be depleted during phytoplankton blooms. See PLANT METABOLISM. Numerous substances are liberated into water by living algae, often with marked ecological effects.
Algae These extracellular products include simple sugars and sugar alcohols, wall polysaccharides, glycolic acid, phenolic substances, and aromatic compounds. Some secreted substances inhibit the growth of other algae and even that of the secreting alga. Some are toxic to fishes and terrestrial animals that drink the water. Nitrogen-fixing blue-green algae release a large proportion of their fixed nitrogen into the water. As part of the plankton tug-of-war, some algae secrete vitamins while others secrete vitamin-inactivating substances. See ALLELOPATHY. Occurrence. Algae are predominantly aquatic, inhabiting fresh, brackish, and marine waters without respect to size or degree of permanence of the habitat. They may be planktonic (free-floating or motile) or benthic (attached). Benthic marine algae are commonly called seaweeds. Substrates include rocks (outcrops, boulders, cobbles, pebbles), plants (including other algae), animals, boat bottoms, piers, debris, and less frequently sand and mud. Some species occur on a wide variety of living organisms, suggesting that the hosts are providing only space. Many species, however, have a restricted range of hosts and have been shown to be (or are suspected of being) at least partially parasitic. Certain algae parasitize fishes while various unicells are characteristic symbionts of protozoa, coelenterates, mollusks, and flatworms. All reef-building corals contain dinoflagellates, without which their calcification ability is greatly reduced. Different phases in a life history may have different substrate preferences. The depth at which algae can grow is determined by the availability of light (water transparency). In the ocean, transparency is affected by latitude (angle of incident sunlight), distance from terrestrial runoff, and physical nature of the bottom. In exceptionally clear areas, such as the Mediterranean and off the east coast of Florida, algae grow at depths of at least 330 ft (100 m). See PHYTOPLANKTON. Many fresh-water algae have become adapted to a nonaquatic habitat, living on moist soil, masonry and wooden structures, and trees. A few parasitize higher plants (expecially in the tropics), producing diseases in such crops as tea, coffee, and citrus. There are several specialized habitats. In deserts, algae (chiefly blue-greens) live on the undersurface and in cracks of rocks. Representatives of several classes (mostly unicells) grow on snow and ice, which they tint green or red depending upon the amount of protective reddish-orange pigment present in the cells. Thermophilic algae (again, chiefly blue-greens) live in hot springs at temperatures up to 163◦F (73◦C), forming a calcareous deposit known as tufa. See TUFA. One of the most remarkable adaptations of certain algae (blue-greens and greens) is their coevolution with fungi to form a compound organism, the lichen. Most lichens are subaerial, being especially common on the bark and dead twigs of trees and on rocks exposed to moisture-laden winds, but a few are marine. See LICHENS. Geographic distribution. Fresh-water algae, which are distributed by spores or fragments borne by the wind or by birds, tend to be widespread if not cos-
mopolitan, their distribution being limited by the availability of suitable habitats. Certain species, however, are characteristic of one or another general climatic zone, such as cold-temperate regions or the tropics. Marine algae, which are spread chiefly by water-borne propagules or reproductive cells, often have distinctive geographic patterns. Along temperate shores, individual species tend to have fairly limited ranges, the end points being determined primarily by water temperature. Thus, the marine floras of Japan, Pacific North America, South Africa, southern Australia, and New Zealand all have a high proportion of endemic species. Distributional patterns of marine algae are correlated with patterns of general oceanic circulation, the main currents not only providing a means and determining the direction of transport of propagules but also controlling the temperature. Major floristic provinces include the North Pacific (Japan eastward through Baja California, Mexico, established by the Kuroshio Current System), the North Atlantic (Cape Hatteras, North Carolina, northeastward through Atlantic Europe, established by the Gulf Stream System), the tropical Atlantic (established by the Atlantic Equatorial Current System), the Arctic (related to the Arctic Ocean), the Antarctic/Subantarctic (related to the Antarctic Convergence and the circumpolar current), and the IndoPacific (transcending modern current systems and related to ancient positions of land masses). The marine algal flora at any given latitude on the Pacific coast of America is at least 90% different from that at the same latitude on the Atlantic coast. See POPULATION DISPERSAL; POPULATION DISPERSION. Many taxonomic groups are widely distributed, but others are characteristic of particular climatic zones or geographic areas. The large brown seaweeds called kelps are mostly confined to cold waters of the Northern Hemisphere, although certain genera occur only in the Southern Hemisphere. Sargassum, also a brown alga, is extremely common in warm waters everywhere. Coenocytic green algae, many of which have calcareous thalli, are characteristic of all tropical areas except the eastern Pacific. Crustose coralline algae are common in all oceans, but individual genera have distribution patterns related to climatic zones and hemispheres. There is a large pantropical floristic element comprising red, brown, and green seaweeds. See PLANT GEOGRAPHY. Vertical distribution. On coasts that are subjected to tidal fluctuation and offer the protection of regular fog or cloud cover, an intertidal belt of seaweeds may be found. When this belt is at least 3 ft (1 m) in breadth, it is often subdivided into horizontal bands that are determined by the ability of particular species to withstand desiccation and wave exposure and by the availability of suitable substrate. Intertidal zonation of seaweeds is correlated with that of invertebrates. Much less is known about patterns of vertical subtidal zonation, where the determining factors appear to be water movement (surge and currents), light intensity, substrate, and predation. Thus, some species grow only at relatively great depths or
435
436
Algae midge larvae
filamentous algae
diatom plankton
marine diatoms and dinoflagellates
copepods
crustacea
small fishes
gizzard shad
plankton algae
ALGAE
black bass
game fish HUMANS
tadpoles plant detritus
carp hake
Chara marine copepods
herring clams oysters
Fig. 2. Food chains from algae to humans. (After H. C. Sampson et al., eds., Textbook of Botany, 3d ed., Harper, 1966)
only on vertical rock faces or only on mollusk shells on quiet sandy bottoms. Algae growing on shores of lakes with wave action (such as Lake Michigan) exhibit small-scale vertical zonation. Economic importance. Planktonic algae, as the primary producers in oceans and lakes, support the entire aquatic trophic pyramid and thus are the basis of the fisheries industry (Fig. 2). Concomitantly, their production of oxygen counteracts its uptake in animal respiration. The ability of certain planktonic algae to assimilate organic nutrients makes them important in the treatment of sewage. See FOOD WEB; SEWAGE TREATMENT. Numerous red, brown, and green seaweeds as well as a few species of fresh-water algae are consumed by the peoples of eastern Asia (especially Japan), Indonesia, Polynesia, and the North Atlantic. Porphyra (nori), Laminaria (kombu), and Monostroma (aonori) are produced by intensive mariculture in Japan and China. Large brown seaweeds may be chopped and added to poultry and livestock feed or applied whole as fertilizer for crop plants. The purified cell-wall polysaccharides of brown and red algae (alginate, agar, carrageenan) are used as gelling, suspending, and emulsifying agents in numerous industries, including food processing, cosmetics, pharmaceuticals, textiles, paper, and printing, as well as in dentistry, medicine, and biological research. Some seaweeds have specific medicinal properties, such as effectiveness against worms. The blue-green alga Spirulina, a traditional food in parts of Mexico and central Africa, is being grown commercially and marketed as a high-protein dietary supplement. Nitrogen-fixing blue-green algae growing in rice paddies provide significant nitrogen enrichment. Petroleum is generally believed to result from bacterial degradation of organic matter derived primarily from planktonic algae. Diatomaceous earth and diatomite (fossil diatom deposits) have many industrial appli-
cations. Unicellular algae, because of the ease with which they can be grown in pure culture and thus provide quantities sufficient for experimentation, have been used traditionally in certain biochemical, biophysical, and physiological studies. Chlamydomonas, Chlorella, Ankistrodesmus, and Euglena have thus been the subject of thousands of investigations, yielding fundamental biological information. See AGAR; ALGINATE. On the negative side, algae can be a nuisance by imparting tastes and odors to drinking water, clogging filters, and making swimming pools, lakes, and beaches unattractive. Sudden growths (blooms) of planktonic algae can produce toxins of varying potency. In small bodies of fresh water, the toxin (usually from blue-green algae) can kill fishes and livestock that drink the water. In the ocean, toxins produced by dinoflagellate blooms (red tides) can kill fishes and render shellfish poisonous to humans. Ciguatera poisoning in tropical areas has been traced to toxin produced by benthic dinoflagellates and accumulated in reef-dwelling fishes. Aquarium fishes may die from an infection by a parasitic dinoflagellate, while humans may suffer from protothecosis, a gastrointestinal disease caused by a colorless chlorophycean unicell (Prototheca). In warm humid regions, leaves of crop plants (tea, coffee, citrus) may be badly damaged by a parasitic green alga. A few marine algae have spread far beyond their original home and become noxious weeds, damaging shellfish plantings, clogging the propellers of small boats, and contributing a large biomass to odoriferous beach wrack. See EUTROPHICATION; TOXIN. Fossil algae. At least half of the classes of algae are represented in the fossil record, usually abundantly, in the form of siliceous, calcareous, or organic remains, impressions, or indications. Blue-green algae were among the first inhabitants of the Earth, appearing in rocks at least as old as 2.3 billion years.
Algebra Their predominance in shallow Precambrian seas is indicated by the extensive development of stromatolites, laminated reeflike structures produced by algal mats that entrap detrital sediments and sometimes deposit calcium carbonate. Modern stromatolites (as at Shark Bay, Western Australia) implicate blue-green algae, although bacteria may have been at least partly responsible for some of the ancient formations. See STROMATOLITE. All three classes of seaweeds (reds, browns, and greens) were well established by the close of the Precambrian, 600 million years ago (mya). Shallowwater reef deposits were formed in various geological periods by calcareous red and green seaweeds, and the process continues today. The chlorophycean order Dasycladales, in which the entire thallus may be calcified, has left a rich record of more than 150 genera dating from the Precambrian. The calcified zygospores (gyrogonites) of Charophyceae are found in fresh-water deposits as old as the Devonian (400 mya). The relatively low diversity among living dasyclads and charophytes clearly indicates that they are relicts. Among red algae, coralline algae date from the Devonian, while two frequently encountered groups (Solenoporaceae and Gymnocodiaceae) are known only as fossils. See GEOLOGIC TIME SCALE. By far the greatest number of fossil taxa belong to classes whose members are wholly or in large part planktonic. Siliceous frustules of diatoms and endoskeletons of silicoflagellates, calcareous scales of coccolithophorids, and highly resistant organic cysts of dinoflagellates contribute slowly but steadily to sediments blanketing ocean floors, as they have for tens of millions of years. Cores obtained in the Deep Sea Drilling Project have revealed an astounding chronology of the appearance, rise, decline, and extinction of a succession of species and genera. From this chronology, much can be deduced about the climate, hydrography, and ecology of particular geological periods. More than 350 genera of dinoflagellates and an even greater number of genera of coccolithophorids (Prymnesiophyceae) have been described from strata younger than the Permian (about 250 mya). Centric diatoms appear in the Cretaceous (about 100 mya), with pennate diatoms apparently evolving in the Paleocene (65 mya). Because of their abundance, planktonic algae have been important rock builders in shallow seas. Mesozoic coccolithophorids resulted in the accumulation of extensive chalks and limestones (such as the white cliffs of Dover) and there are massive Tertiary deposits of diatoms (diatomaceous earth or diatomite). See PALEOBOTANY. Paul C. Silva; Richard L. Moe Bibliography. H. C. Bold and M. J. Wynne, Introduction to the Algae: Structure and Reproduction, 1978; V. J. Chapman, Seaweeds and Their Uses, 1970; J. D. Dodge, The Fine Structure of Algal Cells, 1973; E. Gantt (ed.), Handbook of Phycological Methods: Developmental and Cytological Methods, 1980; J. A. Hellebust and J. S. Craigie (ed.), Handbook of Phycological Methods: Physiological and
Biochemical Methods, 1978; C. S. Lobban and W. J. Wynne (ed.), The Biology of Seaweeds, 1981; F. E. Round, The Ecology of Algae, 1981; J. R. Stein (ed.), Handbook of Phycological Methods: Culture Methods and Growth Measurements, 1973; W. D. P. Stewart (ed.), Algal Physiology and Biochemistry, 1974; H. Tappan, The Paleobiology of Plant Protists, 1980; O. R. Zaborsky (ed.), CRC Handbook of Biosolar Resources, vol. 1, pt. 1, 1982.
Algebra The branch of mathematics dealing with the solution of equations. These equations involve unknowns, or variables, along with fixed numbers from a specified system. The origins of algebra were based on the need to develop equations that modeled real-world problems. From this came a very extensive theory based on the need to find the values that can be successfully used in the equations. Number systems. Classical algebra is conducted in one of several number systems. The most basic is the natural numbers, consisting of the counting numbers: 1, 2, 3, 4. . . . The natural numbers along with their negatives and 0 form the set of integers: . . . , −2, −1, 0, 1, 2,. . . . Any number which can be expressed as a quotient a/b of two integers a and b is called a rational number. If the decimal expansion of a number either is repeating or terminates in a string of zeros, it is always possible to find integers a and b such that the quotient a/b gives the original number. Numbers where the decimal expansion cannot be written as a quotient of integers are the so-called irrational numbers. Early mathematicians did not recognize the existence of such numbers until geometric considerations showed that they must exist as the lengths of sides of right triangles. The sets of rational numbers and irrational numbers have no elements in common, and together they make up the real-number system. It is within the system of real numbers that algebraic tasks are most often performed. However, there is a still larger set of numbers, known as the complex numbers, which serve much more completely in solving equations. The complex numbers consist of numbers √ of the form a + bi, where i is a symbol representing − 1. The symbol i is commonly called the imaginary root of −1, and the numbers of the form bi are known as the set of imaginary numbers. In a complex number a + bi, the real number a is called the real part while the real number b is the imaginary part. See COMPLEX NUMBERS AND COMPLEX VARIABLES. Real numbers are often represented on a number line (Fig. 1a). In this method of depiction, there is a one-to-one correspondence between the points on the line and the real numbers. There is no such representation of the complex numbers since the number line implies an ordering notion which does not easily extend beyond the real numbers. However, the complex numbers can be represented on a
437
438
Algebra
–2
0
–1
1
2
(a)
imaginary axis
a
real axis
The symbols are usually called variables, while the numbers (or constants) multiplied by them are the coefficients of the expression. The coefficients are all assumed to come from a designated number system. Algebraic expressions take on specific values when each of the variables involved is assigned a numerical value. Algebraic equations. A statement is a relationship between several algebraic expressions. This relationship can be equality (=) or an inequality (< or >), but in any case it puts restrictions on the values of the variables. For example, each of expressions (2) x2 + 3xy = y − 3
b
2xy > y5 − 1
represents a relationship between the variables x and y, while Eq. (3) relates the three variables, x, y, and z.
a + bi
4z − 2xy = 5 + yz2
(b) Fig. 1. Representations of number systems. (a) Of real numbers on a number line. (b) Of complex numbers on a two-dimensional plane.
two-dimensional plane (Fig. 1b), by plotting the real part on the horizontal axis and the imaginary part on the vertical axis. Operations. Within any number system there is a set of operations which can be performed on some or all of the members of the system. Operations take a pair of numbers in the system and produce a single new number from this pair. The most standard operations are addition (+), multiplication (·), subtraction (−), and division (÷). A set is closed under an operation if the application of the operation to any pair of numbers in the set results in another number in the same set. For example, the set of natural numbers is closed under + but not under − (since, for example, 1 − 2 = −1 is not a natural number). The rational numbers, the real numbers, and the complex numbers are closed under all four of the basic operations (except in the case where division by 0 may occur). See ARITHMETIC. Other operations arise from the four basic ones. Exponentiation is derived from the process of repeatedly multiplying a number by itself. If n is a positive integer, the symbol an indicates that the number a is multiplied by itself n times. If n > 0 is an integer, then a−n refers to the quotient 1/an. Fractional exponents can also be considered by defining a1/n to be the number b such that bn = a. This notion is further extended by letting am /n = (a1/n)m. For real numbers, some fractional exponents may not exist. For example, (−1)1/2 is not a real number. The complex numbers, however, are closed with respect to exponentiation. Algebraic expressions. In algebra, symbols are frequently used to designate unknown values. For example, letters such as x and y can stand for any one of a set of numbers satisfying certain conditions. These symbols are combined with numbers and the basic operations to form algebraic expressions such as (1). 3x2 + 5y3 x − 6
(x5 − 1)3xy y2 + 2 x
(2)
(1)
(3)
If, when specific values from an appropriate number system are substituted for the variables in an equation, the resulting statement is true, then these numbers are said to be a solution to the equation. The solutions of a particular equation are dependent on the number system in use. For instance, the equation x4 = 2 has no solutions in the set of rational numbers, two solutions (x = 21/4 and x = −21/4) in the real numbers, and four solutions in the complex field. Algebraic equations are used to describe many phenomena in the real world. Their application runs through almost all disciplines in the sciences, social sciences, and business. Once an equation is established as a good model of some situation, it is necessary to determine the solutions to the equation, as these represent instances when the conditions described by the equation are actually met. One of the most basic kinds of equations is the linear equation. If such an equation involves the variables x0, x1, . . . , xn, then it can be written in the form of Eq. (4), where a0, a1, . . . , an, b are coefficients a0 x0 + a1 x1 + · · · + an xn = b
(4)
from some number system. In the simplest form, where only two variables, x and y, are involved, a linear equation represents a straight line in two dimensions. The equation of such a line can be written in the form y = mx + b, where m is the slope of the line (the rate of change of the y variable with respect to the x variable) and b is the y intercept, or the y value corresponding to x = 0. Another important algebraic equation is the quadratic equation, which has the standard form y = ax2 + bx + c for coefficients a, b, c. The values of x that produce the value y = 0 are found through the quadratic formula, Eq. (5). These x values are known √ −b ± b2 − 4ac x= (5) 2a as roots. In the complex number system, there are always either one or two roots to a quadratic equation. However, in the real number system, the number of roots of a particular quadratic equation is
Algebra given by the discriminant, b2 − 4ac. There are two real roots to a quadratic equation if the discriminant is positive, one real root if it is 0, and no real roots if it is negative. For example, the equation y = 3x2 − 2x + 1 has the discriminant given by Eq. (6), and b2 − 4ac = (−2)2 − 4(3)(1) = −8
y axis
(a, b)
b
a
(6)
hence has no real roots. The complex roots are x1 and x2 given by Eqs. (7). The ability to always find √ √ −(−2) + (−2)2 − 4(3)(1) 1 + 2i x1 = = (7a) 2(3) 3 √ √ 1 − 2i −(−2) − (−2)2 − 4(3)(1) x2 = = (7b) 2(3) 3 roots of quadratic equations (as well as much more complicated ones) in the complex numbers is an indicator of the usefulness of this number system. Functions. Some of the most fundamental algebraic equations are those involving two variables. When one of these variables, say y, can be expressed in terms of the other variable, say x, in such a way that each value of x produces exactly one value of y, then the relationship described by the equation is known as a function. More general relationships which may or may not produce unique values of y from each value of x are called relations. In a function such as has been described, the variable x is called the independent variable while the resulting variable y is called the dependent variable. The possible values that can be used by the independent variable make up the domain of the function (or, more generally, the relation). The resulting values of the dependent variable compose the range of the function. It is most common and quite practical to write a variable y which is dependent on a variable x in the form y = f(x). This notation, known as function notation, indicates that every y value is determined by an x value. An equation such as y = x2 + 3 would be written in functional notation as f(x) = x2 + 3. If x = a, then the corresponding value of y is given by y = f(a). In the particular function f(x) = x2 + 3, the value of y associated with x = −1 is y = f(−1) = (−1)2 + 3 = 4. Suppose an equation involves the variables x and y. It is said that the ordered pair (a, b) is part of the function (or relation) relating x and y whenever the value x = a produces the value y = b. The ordered pair (−1, 4) is part of the function f(x) = x2 + 3, since f(−1) = 4. Equations involving three variables, say x, y, z, produce ordered triples (a, b, c) describing corresponding values of the variables. This process can readily be extended so that an equation with n variables would give rise to a set of ordered “n-tuples” where the terms for the corresponding variables can be substituted successfully into the equation. Graphs. Relations, and more specifically functions, are frequently represented effectively with graphs. If an equation involves the variables x and y, then a two-dimensional coordinate system can be used to depict the values that correspond in the equation. If
x axis
(a)
y
x
(b) Fig. 2. Representation of relations by graphs. (a) Graphic depiction of an ordered pair (a, b) given by an equation. (b) Graph of the function f(x) = x2 + 3.
the ordered pair (a, b) is given by the equation, then it can be graphically depicted (Fig. 2a). All of the ordered pairs given by the equation combine to give a complete graphic representation of the relationship described by the equation (Fig. 2b). Functions are readily distinguished among more general relations by their graphical representations. Since each x value from the domain of a function produces only one y value in the range, and since all of the points on a vertical line have the same x value, it must be the case that each vertical line intersects the graph of a function in at most one place. For example, from the vertical line test it is readily seen that the relation described in Eq. (8) is a function y = 2x2 − 4x + 5
(8)
(Fig. 3a) while the relation described by Eq. (9) is not (Fig. 3b). y2 = x − 1
(9)
One of the most important properties of a function is given by the places where the dependent variable takes on the values 0. These correspond to the points where the graph crosses, or touches, the horizontal or x axis. The values of x for which f(x) = 0 are called zeros (or roots) of the function. See ANALYTIC GEOMETRY; GRAPHIC METHODS. Polynomials. There are many special functions which are commonly used to describe real situations. One of the most basic is the polynomial, which is a function of the form (10). Each coefficient a0, f (x) = a0 + a1 x + a2 x2 + · · · + an xn
(10)
439
440
Algebra plex number polynomial f(x) can be factored completely as f(x) = (x − a1) (x − a2) . . . (x − an), where a1, a2,. . . , an are the complex number roots and n is the degree of the polynomial. This is not always possible for real polynomials. For example, the polynomial f(x) = x2 + 1 cannot be factored further by using polynomials with real coefficients. It is called irreducible over the real numbers. As a complex polynomial, it can be factored as x2 + 1 = (x + i ) · (x − i ). Factoring of polynomials. Even in the real numbers, a polynomial often may be factored into other polynomials of lesser degree. This allows for analysis of the polynomial by consideration of the simpler factors. Some of the most common and useful factoring patterns are given in Eqs. (12).
y (0,5)
x
(a)
y
(1,0)
x
x2 − a2 = (x − a)(x + a)
(12a)
x2 + 2ax + a2 = (x + a)(x + a)
(12b)
x2 − 2ax + a2 = (x − a)(x − a)
(12c)
x3 + 3ax2 + 3a2 x + a3 = (x + a)(x + a)(x + a) x + 3ax − 3a x − a 3
2
2
= (x − a)(x − a)(x − a)
(b) Fig. 3. Graphs of the relations described by the equations (a) y = 2x2 − 4x + 5 and (b) y2 = x − 1. The vertical-line test shows that the first relation is a function while the second is not.
a1, . . . , an is a real number, and the exponents of the variables are nonnegative integers. The largest exponent of x is called the degree of the polynomial. When graphed, a polynomial of degree n will reverse directions no more than n − 1 times. For example, the graph of the degree-4 polynomial in Eq. (11) [Fig. 4] changes direction three times.
(12d )
3
(12e)
x3 − a3 = (x − a)(x2 + ax + a2 )
(12f )
x + a = (x + a)(x − ax + a )
(12g)
3
3
2
2
Rational functions. Functions which are of the form f(x) = g(x)/h(x), where g(x) and h(x) are polynomials, are called rational functions. They frequently appear in instances where comparisons of two polynomials are necessary. A rational function is defined only at points where the denominator is nonzero. The roots of a rational function agree with y
f (x) = 3x4 − 4x3 − 12x2 + 8
(11)
A polynomial of degree n will have no more than n roots in the real number system. Sometimes it may have fewer. For example, the function f(x) = x2 + 1 has no real roots, while the function f(x) = x2 − 1 has the roots x = − 1 and x = 1. The function f(x) = x2 has the single root x=0. Whenever x = a is a root of a polynomial f(x) of degree n, then the polynomial can be factored as f(x) = (x − a) g (x), where g (x) is a polynomial of degree n − 1. If the term (x − a) can be successively factored from a polynomial k times, then x = a is said to be a root of multiplicity k. The equation y = (x −1)(x −1) has only one root, x = 1, but this root is of multiplicity 2. It is in the consideration of roots of an equation that the complex number system becomes particularly important. If a polynomial is formed with coefficients from the complex number system, then the sum of the multiplicities of the roots (from the complex numbers) of the polynomial coincides with the degree of the polynomial. Consequently, a com-
(0, 8)
(–1 , 3)
x
(2 , – 24) Fig. 4. Graph of the degree-4 polynomial f(x) = 3x4 − 4x3 − 12x 2 + 8, which changes direction three times.
Algebra y
applied to b produces x, as in Eqs. (15). The logay = logb x ↔ x = by
rithm function, like the exponential function, has a characteristic graph (Fig. 6). Basic rules that are necessary for calculation with logarithms are given in Eqs. (16). For example, since
(0,1)
x
(a)
y
(0,1)
x
(b) Fig. 5. Graphs of exponential functions. (a) Decreasing graph of y = ( 12 )x. (b) Increasing graph of y = 2x.
the roots of the polynomial that is found in the numerator. Exponential functions. In general, an exponential function is one which has the form f(x) = bx for some constant b in the number system. The graphs of these functions take on one of two basic types. The graph is decreasing if 0 < b < 1 (Fig. 5a), and it is increasing if b > 1 (Fig. 5b). The case of b = 1 is trivial, since all values of the function are 1. Exponential functions are of special importance in describing a situation where the rate of growth is proportional to the amount present. Among the situations that can be described by exponential functions are those involving the calculation of interest, radioactive decay, and population growth. Calculation with exponential functions is based on the basic equalities given in Eqs. (13). bx+y = bx by bx−y =
(13a)
x
b by
(15)
(13b)
(ab)x = ax bx
(13c)
(bx )y = bxy
(13d)
logb (xy) = logb x + logb y x = logb x − logb y logb y
(16a)
logb xr = r(logb x)
(16c)
logb bx = x
(16d)
blogb x = x
(16e)
(16b)
logarithms are exponents, the exponent for the product in Eq. (16a) is the sum of the exponents for the factors. See LOGARITHM. Systems of linear equations. Frequently, a situation is described by a system of linear equations which must all be satisfied simultaneously. For example, the system of Eqs. (17) involves three variables and three equations. 2x + 3y − 4z = 2 4x − 5y − z = 0 −3x + y − 2z = 1
(17)
More general systems involve n variables and m equations and take the form of Eqs. (18). A solution a11 x1 + a12 x2 + · · · + a1n xn = b1 a21 x1 + a22 x2 + · · · + a2n xn = b2 .. . am1 x1 + am2 x2 + · · · + amn xn = bm
to the system is a set of numbers ci, c2, . . . , cn which satisfies each equation. The equations in a linear system can be multiplied by constants, where a single number is multiplied by each term in the sum. They can also be added together, where the coefficients of the corresponding variables are summed. A linear combination of one or more equations in the system is formed by multiplying each equation by a constant and then y
Logarithmic functions. Another function that plays an especially important role in algebra is the logarithm. In general, if b > 1, the logarithm to the base b of a value x is denoted by Eq. (14). The logarithm y = logb x
(1,0)
x
(14)
function is defined in terms of exponentials by reversing the roles of the x and y variables. It follows that if y = logb, then y is the exponent which when
(18)
Fig. 6. Graph of the logarithm function y = logbx.
441
442
Algebraic geometry summing the resulting equations. When one equation in a system is a linear combination of one or more of the other equations, the system is called dependent. The solutions to a dependent system are the same as the solutions to the system formed by deleting the equation that is a linear combination of other equations in the system. In order to find solutions to a system, it is frequently expedient to first try to eliminate equations that are linear combinations of others in the system. See LINEAR SYSTEMS OF Wayne B. Powell EQUATIONS. Bibliography. D. T. Christy, College Algebra, 2d ed., 1993; C. H. Edwards and D. Penney, Calculus and Analytic Geometry, 4th ed., 1994.
Algebraic geometry The study of zero sets of polynomial equations. Examples are the parabola y − x2 = 0, thought of as sitting in the (x, y)-plane, and the locus of all points (t, t2, t3, t4), which is defined by Eqs. (1), in the y2 = xz
yx = wz
wy = x2
(1)
coordinate space with coordinates w, x, y, z. Another interesting example is an elliptic curve, typically defined by an equation like Eq. (2), where A and B are y2 = x(x − A)(x − B)
(2)
constants. Objects such as these are called affine algebraic sets. They exist in affine n-space, denoted An, which is defined to be the coordinate space with coordinates x1, . . . , xn. The coordinate ring k[V] of an affine algebraic set V is the set of functions on V obtained by restriction from polynomial functions on the ambient affine space. These functions can be added, subtracted, and multiplied (that is, they form a ring). For example, the coordinate ring of the parabola has in it functions y and x satisfying y = x2. This relation can be used to eliminate all references to y, so that the coordinate ring of the parabola is identified with the ring of polynomials in x. See ANALYTIC GEOMETRY; POLYNOMIAL SYSTEMS OF EQUATIONS; RING THEORY. The possibility of studying a question geometrically via the zero loci of polynomials or algebraically via the algebra of the coordinate ring gives the subject much of its power and flavor. This has led to an amazing growth in applications to other disciplines. For example, the integers . . . , −2, −1, 0, 1, 2, . . . form a ring that is algebraically similar to the coordinate ring of the affine line, and algebrogeometric methods have come to play a central role in number theory. In studying the path space of strings and the partition function, modern physicists have made the moduli space of curves a central object of research. Finally, a number of differential equations in engineering and physics can best be studied algebrogeometrically, although their solutions are not algebraic functions. See DIFFERENTIAL EQUATION; NUMBER THEORY; SUPERSTRING THEORY.
Basic numerical invariants. An algebraic set V is irreducible if fg = 0 in k [V] implies that either f is 0 or g is 0. For example, the algebraic set defined by xy = 0 is not irreducible. An irreducible algebraic set is called an algebraic variety. Dimension. The procedure of defining the dimension of an algebraic variety V begins with the intuition that affine n-space has dimension n. Geometrically, if V is defined by the vanishing of several polynomials in n variables, so that V is contained in affine n-space, the coordinates on An are changed in what is said to be a general way and the projection given by Eqs. (3) is considered. It turns out that for An → Ad
(x1 , . . . , xn ) → (x1 , . . . , xd )
(3)
one value of d ≤ n, this mapping when restricted to V will be surjective with finite fibers, that is, there will exist at least one and at most a finite number of points of V of the form (a1, . . . , ad, . . . ) for given a1, . . . , ad. This value of d is the dimension of V. Algebraically, if k[V] denotes the coordinate ring of V, the function field k(V ) is considered, consisting of fractions f/g with f and g in k[V] and g not zero. [The word “field” refers to the fact that in k(V ) it is possible to add, subtract, multiply, and divide.] There will exist, in k(V ), d elements z1, . . . , zd satisfying no polynomial relations among themselves, but having the property that, for any w in k(V ), there is a nonzero polynomial P(x1, . . . , xd+1) with P(z1, . . . , zd, w) = 0. See FIELD THEORY (MATHEMATICS). Degree. The degree of an algebraic set is a central notion of enormous depth and power. Intuitively, if V is a variety of dimension d in An, the degree of V is the number of points of intersection of V with An−d. There are a number of pitfalls here. So far, there has been no specification of where the functions take values. If the coordinates are real numbers, then the parabola y − x2 = 0 meets the line y = 1 in two points, but it meets the line y = −1 in zero points (Fig. 1). This problem is solved by taking coordinates in the field C of complex numbers. (More generally, it is possible to work over any algebraically closed to ∞
y = x2
y=1 x axis y = −1
y axis x = 1 Fig. 1. Parabola in affine 2-space, A2. While it would seem that lines in A2 can meet the parabola in 0, 1, or 2 points, with the introduction of algebraically closed fields of coordinates and projective varieties, all such intersections have two points, counting multiplicity.
Algebraic geometry field. A field F is said to be algebraically closed if any polynomial with coefficients in F has a root in F. The complex numbers have this property.) It can also be argued that the parabola meets the line y = 0 in just one point, but this is a point of tangency and therefore must be counted twice. See COMPLEX NUMBERS AND COMPLEX VARIABLES. A more serious difficulty is how to arrange for the parabola to meet the lines x = constant in two points. A glance at Fig. 1 suggests that one point of intersection should be thought of as lying at ∞ (infinity). The affine (x, y) plane containing the parabola can be imagined positioned as the locus Z = 1 inside A3 with coordinates X, Y, Z. The set of all lines through the origin in A3 meeting the parabola (Fig. 2) is contained in the locus X2 − YZ = 0, but there is one line, defined by Z = X = 0, lying in this locus but not meeting the parabola. The lines through the origin meeting the line x = c lie in the plane X − cZ = 0 in A3, and the common zeros, satisfying Eq. (4), X − cZ = 0 = X 2 − Y Z
(4)
consist of two lines, one being the line that joins (0,0,0) with (c, c2, 1) and the other being the line Z = X = 0. This motivates defining projective n-space, Pn, to be the set of all lines through the origin in An+1. For example, the Riemann sphere P1 is equivalent to A1 together with the point ∞ at infinity. Projective algebraic sets in Pn are defined by the vanishing of homogeneous polynomials in the coordinates on An+1. For many purposes, it is preferable to work with projective algebraic sets, which can be thought of as completions or compactifications of the affine algebraic sets defined above. In the above example, the collection of lines through the origin lying on the homogeneous conic X2 − YZ = 0, which comprises the lines that meet the parabola x2 = y together with the line X = Z = 0, is the projective completion of the parabola. If V is a projective algebraic variety of dimension n in PN and L is a linear space of dimension N − n such that the intersection of V and L is finite, then the number of points of intersection of V and L, counted with multiplicity, is independent of L and is called the degree of V. See PROJECTIVE GEOMETRY. Classification. The problem of classifying projective algebraic varieties typically leads to a finite Z axis
xis
Xa
Z=0
Y axis Z=X=0
(b)
(c)
(0,0,0)
Fig. 2. Projective completion of the parabola of Fig. 1, consisting of lines through the origin lying on the homogeneous conic X2 − YZ = 0, including lines meeting the parabola x2 = y together with the line Z = X = 0.
Ioop ms
Fig. 3. Riemann surfaces of low genus. (a) Genus 0. (b) Genus 1. (c) Genus 2.
number of discrete invariants together with a continuous parameter space, the moduli space of the problem, which is itself an algebraic variety. The classification of projective varieties of dimension 1 (algebraic curves) and dimension 2 (algebraic surfaces) will be considered. Algebraic curves. To help visualize the classification of algebraic curves, the coordinates will be chosen to be complex numbers, and the curves will be assumed to be nonsingular, that is, to have a welldefined tangent at each point. In this case, the algebraic curves are topological surfaces (frequently referred to as Riemann surfaces), and, by topology, they can all be realized as (hollow) doughnuts with g = 0, 1, 2, . . . holes (Fig. 3). The number of holes is the genus of the curve and is the only discrete invariant. Algebraically, a polynomial in two variables P(x, y) can be written as in Eq. (5). After projective complea0 (x)yn + a1 (x)yn−1 + · · · + an (x) = 0
(5)
tion and elimination of singular points (by a process called normalization), the locus of zeros of P yields an algebraic curve C. For a general value of x, there are n solutions for y, and therefore C can be visualized as an n-sheeted covering of the projective line P1. The function y has poles on C, but it is single-valued and satisfies the given polynomial relation. Simple, but by no means trivial, examples are the elliptic and hyperelliptic equation (6). Values of x (possibly including y2 = f (x)
(6)
x = ∞) for which there are fewer than n solutions in y are called branch points. For a general polynomial P(x, y) there will be n − 1 solutions in y over branch points. In this case, the number b of branch points is even, and the genus of C is given by the Hurwitz formula (7). g = 1 − n + (b/2)
Z=1
x2 = y
(a)
(7)
See CONFORMAL MAPPING; TOPOLOGY. Curves of a given genus form a continuous family parametrized by the moduli space Mg. The only curve of genus zero is the Riemann sphere, and therefore Mg is a point. Curves of genus 1 (elliptic curves) are complex tori, which can be viewed as quotients of the complex plane modulo a lattice Lz (called the period lattice) consisting of elements of the form m + nz for fixed complex number z and all integers m and n. Scaling the lattice by a complex factor does not change the quotient curve, and therefore, if a, b, c, and d are integers satisfying Eq. (8) and
443
444
Algebraic geometry w is given by Eq. (9), then Lw and Lz give the same curve. [Indeed, Lz is related to Lw by Eq. (10).] Alad − bc = ±1 w = (az + b)(cz + d)
and H2(V, L). These are finite-dimensional vector spaces. The Riemann-Roch theorem calculates the Euler-Poincar´e characteristic given by Eq. (12).
(8) −1
Lz = (cz + d)Lw
(9) (10)
gebraic functions on M1 are √ meromorphic functions f(z) defined for z = x + y − 1 and y = 0, satisfying f(z) = f(w) with any w as above, and not growing too rapidly as y → ∞. These coincide with rational functions P[ j(z)]/Q[ j(z)] in the classical j-function given by Eq. (11). In fact, M1 = A1 with parameter j(z) = q−1 (1 + 744q + 196884q2 + · · ·)
(11)
q = e2π iz j. For g > 1, Mg is an algebraic variety of dimension 3g − 3. Although not complete, it may be completed to a projective variety by admitting points corresponding to certain singular curves (stable curves). A divisor D on a curve C is a formal sum n1p1 + n2 p2 + ··· + nrpr with ni an integer and each pi a point of C. The degree of D equals n1 + n2 + ··· + nr. To an algebraic function f on C is associated a divisor (f ) of degree 0 by associating to each singular point its multiplicity as a zero or pole. By definition, the jacobian J(C ) is the group of divisors of degree 0 modulo the subgroup of divisors (f ) associated to functions. The jacobian J(C ) is a g-dimensional complex torus, whose period lattice is given by integrating a basis of g algebraic differential one-forms over closed loops. In fact, J(C ) is a projective algebraic variety (called a principally polarized abelian variety), and the classical Torelli theorem states that the association C → J(C ) embeds Mg into the Siegel moduli space Ag of principally polarized abelian varieties of dimension g. The Schottky problem of identifying the image of Mg in Ag has been solved. Algebraic surfaces. The classification of projective varieties of dimension 2 (algebraic surfaces) hinges on a careful analysis of line bundles and linear systems. To map a variety V to Pn means to associate to each point of V a line through the origin in An+1, but a modern algebraic geometer would take a dual view and think of the space of linear functions on the line. The collection of these dual lines associated to points of V forms a line bundle. A linear function f on An+1 induces an element fυ in the line over υ in V, and hence a section of the line bundle. The set of υ where the fυ vanishes is a Cartier divisor on V. A twodimensional subspace of linear functions gives rise to a linearly varying family of Cartier divisors called a pencil on V. This correspondence between line bundles, linear systems, and Cartier divisors lies at the heart of the classification of algebraic surfaces. The key difficulty in calculating the dimension of the space of all sections of a given line bundle is usually overcome via the Riemann-Roch theorem. The space of sections of line bundle L is denoted H 0(V, L), but there are also defined higher cohomology groups H1(V, L)
χ(L) = dimension of H 0 (V, L) − dimension of H 1 (V, L) + dimension of H 2 (V, L) (12) Among the interesting kinds of surfaces are (1) rational surfaces, whose function field k(V ) is generated by two independent transcendental elements (P2 is a rational surface); (2) ruled surfaces, with a one-parameter family of curves of genus 0; (3) elliptic surfaces, with a pencil of elliptic curves; and (4) K3 surfaces, a family of surfaces including nonsingular surfaces of degree 4 in P3. These K3 surfaces are of interest to physicists as providing a plausible model for the universe. Applications to differential equations. The linear Picard-Fuchs equations arise very naturally. If [Xs] is a family of algebraic curves (or indeed of varieties of any dimension), depending on a parameter s, then the de Rham cohomology of a curve Xs can be thought of as the space of closed differential forms modulo exact forms, that is, as the space of integrands to be integrated over closed loops ms on Xs (Fig. 3). These integrands can themselves be differentiated with respect to the parameter s. Since they form a finite-dimensional space, there will be some operator having the form given by Eq. (13) D = a0 (s)(d/ds)n + a1 (s)(d/ds)n−1 + · · · + an (s) (13) [Picard-Fuchs equation] which vanishes identically. A solution may be obtained by choosing a de Rham cohomology class ws and a loop ms varying continuously with s. The integral in Eq. (14) satisfies Eq. (15). As s describes a loop in the parameter F(s) = ws (14) DF(s) = Dws = 0 (15) space, ms does not necessarily return to its original position, and therefore F(s) is in general a multiplevalued function of s. For example, the hypergeometric equation (16) can be studied in this way. s(1 − s)(d 2 f/ds2 ) + [c − (a + b + 1)s](d f/ds) −ab f = 0 (16) See HYPERGEOMETRIC FUNCTIONS. There has also been much interest in certain nonlinear partial differential equations such as the Korteweg–de Vries equation (17). The coordinate ∂f/∂t = ∂ 3 f/∂s3 + 6f ∂f/∂s
(17)
ring R of an affine algebraic curve is embedded inside a large ring D of differential operators of the form
Algorithm D = a0(s)(d/ds)n + · · · mentioned above. This embedding depends on the choice of a line bundle L. When L is made to vary with respect to a second parameter t, the Korteweg–de Vries equation and related equations arise as equations of motion for a fixed element of R inside D. Applications to number theory. Three major conjectures have been solved by using algebrogeometric techniques. The Weil conjecture concerned numbers of solutions of equations mod p. A key idea was to reinterpret solutions modulo p as fixed points of the Frobenius endomorphism given by Eq. (18). The (x1 , . . . , xn ) → (x1 , . . . , xnp ) p
(18)
Lefschetz fixed-point theorem in algebraic topology counts the number of fixed points of an endomorphism. This theorem was ported to algebraic geometry via e´tale cohomology and used to help solve the conjecture. The Mordell conjecture stated that a curve of genus greater than 1 can have only a finite number of points whose coordinates are rational numbers. This implies, for example, that many polynomials in two variables, like the Fermat equation (19), have x n + yn − 1 = 0
(n ≥ 4)
(19)
only finitely many rational solutions. The proof of the Mordell conjecture involved an arithmetic generalization of the concept of degree, which was then applied to study the geometry of the Siegel moduli space. Finally, it seems likely that the Fermat conjecture itself, which states that Eq. (19) has no rational solutions with xy = 0, has been proved. Thus it can be said that Eq. (20) has no solution in integers with n an + bn + cn = 0
(20)
prime and abc = 0. Central to the proof is the Frey elliptic curve, given by Eq. (21), which is associated y2 = x(x − an )(x + bn )
(21)
to a solution of Eq. (20). These problems illustrate the power of geometry in dealing with seemingly algebraic problems like the Mordell and Fermat conjectures which admit no plausible, purely algebraic attack. Even the study of mathematics at the school and college level, where the algebraic manipulation of formulas tends to be emphasized, can benefit from greater attention to deeper and more powerful geometric insights. See GEOMETRY. Spencer J. Bloch Bibliography. S. S. Abhyanker, Algebraic Geometry for Scientists and Engineers, 1992; J. Harris, Algebraic Geometry: A First Course, 1992; R. Hartshorne, Algebraic Geometry, rev. ed., 1991; I. R. Shafarevich, Basic Algebraic Geometry, 1977, paper, 1990.
Alginate A major constituent (10–47% dry weight) of the cell walls of brown algae. Extracted for its suspending, emulsifying, and gelling properties, it is one of three algal polysaccharides of major economic importance, the others being agar and carrageenan. See AGAR; CARRAGEENAN. Alginate is a linear (unbranched) polymer in which regions have a predominance of either D-mannuronic acid with β-1,4 linkages or L-guluronic acid with α1,3 linkages. The ratio of the monomers varies greatly among different species, in the same species at different seasons, and in different parts of the same plant. Regions rich in mannuronic acid have a ribbonlike conformation, while those rich in guluronic acid have kinks with a marked affinity for calcium ions. Alginate gels are not formed by cooling, as with agar and carrageenan, but by the addition of calcium ions, which causes separate polymeric chains to cohere along surfaces rich in guluronic acid. The chief sources of alginate are members of the family Fucaceae (rockweeds) and the order Laminariales (kelps), harvested from naturally occurring stands on North Atlantic and North Pacific shores. The supply has been significantly increased through mariculture, especially in China. Processing is carried out chiefly in the United Kingdom, United States, Norway, Canada, France, and China. In the extraction process, washed and macerated plants (fresh or dried) are digested with alkali to solubilize the alginate, and the filtrate is precipitated with a concentrated solution of calcium chloride. The calcium alginate is converted to alginic acid by dilute hydrochloric acid. After further purification and chemical treatment, the product is marketed as sodium, potassium, or ammonium alginate, usually in powdered form. Because of its colloidal properties, alginate finds numerous industrial applications, especially in the food, textile, paper, printing, paint, cosmetics, and pharmaceutical industries. About half of the consumption is in the making of ice cream and other dairy products, in which alginate prevents the formation of coarse ice crystals and provides a smooth texture. As an additive to paint, it keeps the pigment in suspension and minimizes brush marks. An alginate gel is used in making dental impressions. See ICE CREAM; MILK; PHAEOPHYCEAE. Paul C. Silva; Richard L. Moe
Algorithm A well-defined procedure to solve a problem. The study of algorithms is a fundamental area of computer science. In writing a computer program to solve a problem, a programmer expresses in a computer language an algorithm that solves the problem, thereby turning the algorithm into a computer program. See COMPUTER PROGRAMMING. Classic algorithms. Efficient algorithms for solving certain classic problems were discovered long before
445
446
Algorithm the advent of electronic computers. Sieve of Eratosthenes. The algorithm known as the Sieve of Eratosthenes dates back to the 3d century B.C. It solves the problem of finding all the prime numbers between 1 and a positive integer N. It works as follows: First write the numbers from 1 to N, and cross out 1. Then cross out the multiples of 2 except 2 itself. Then find the first number not yet crossed out after 2. The number is 3. Cross out the multiples of 3 except 3 itself. Then find the first number not yet crossed out after 3. The number is 5. Cross out the multiples of 5 except 5 itself. And so on. This process is repeated until no more numbers can be crossed out. For example, the case N = 20, carried through the crossing out of multiples of 3, gives: − − 11 − − 1 2 3− 4 5 − 6 7− 8 − 9 − 10 12 − 19 − − − − − − 17 − − 13 − 14 15 16 18 20 Consideration of the numbers 5, 7, and so on will not cross out any more numbers. The numbers that have not been crossed out (2, 3, 5, 7, 11, 13, 17, 19) are the prime numbers between 1 and 20. See NUMBER THEORY. Euclid’s algorithm. Euclid’s algorithm, detailed in his treatise Elements, solves the problem of finding the greatest common divisor of two positive integers. It makes use of the fact that the greatest common divisor of two numbers is the same as the greatest common divisor of the smaller of the two and their difference. It works as follows: First find the larger of the two numbers and then replace it by the difference between the two numbers. Then find the larger of the resultant two numbers and then replace it by the difference between the two numbers. And so on. When the two numbers become equal, their value is the greatest common divisor and the algorithm terminates (finishes). The steps of the algorithm when it is used to find the greatest common divisor of 48 and 18 are: − Replace 48 by 30 (= 48 − 18) − 48 − Replace 30 by 12 (= 30 − 18) − 30 Replace 18 by 6 (= 18 − 12) 12 − Replace 12 by 6 (= 12 − 6) − 12 The answer is 6 6
18 18 − − 18 6 6
Operation. An algorithm generally takes some input (for example, the value N in the Sieve of Eratosthenes and the two given integers in Euclid’s algorithm), carries out a number of effective steps in a finite amount of time, and produces some output. An effective step is an operation so basic that it is possible, at least in principle, to carry it out using
front A
A1
Ai
A2
Fig. 1. An array of n elements.
top
C
inserting element D
queue
D C
removing element D
B
B
A stack
A
A
Fig. 2. Inserting and deleting an element from a stack.
pen and paper. In computer science theory, a step is considered effective if it is feasible on a Turing machine or any of its equivalents. A Turing machine is a mathematical model of a computer used in an area of study known as computability, which deals with such questions as what tasks can be algorithmically carried out and what cannot. See AUTOMATA THEORY; RECURSIVE FUNCTION. Data structures. Many computer programs deal with a substantial amount of data. In such applications, it is important to organize data in appropriate structures so as to make it easier or faster to process the data. In computer programming, the development of an algorithm and the choice of appropriate data structures are closely intertwined, and a decision regarding one often depends on knowledge of the other. Thus, the study of data structures in computer science usually goes hand in hand with the study of related algorithms. Commonly used elementary data structures include records, arrays, linked lists, stacks, queues, trees, and graphs. A record is a group of related components. For example, a record representing a date may have three components: month, day, and year. An array is a sequence of indexed (subscripted) components, A1, A2, . . . , An, stored in the memory of a computer in such a manner that each component Ai can be accessed by specifying its index (subscript) i (Fig. 1). The components of an array can be accessed individually and independently. In contrast, a linked list is a sequence of components stored in such a manner that the components can be accessed only sequentially. That is, the first component must be accessed before the second, the second before the third, and so on. A stack is a sequence of components that allows component access, addition, and removal only at one end of the sequence called the top of the stack (Fig. 2). Adding an element to the top of a stack is
front
C
A inserting element D
Fig. 3. Inserting and removing an element from a queue.
B
C
B
rear B
An
C
D
B removing element A
rear
C
D
Algorithm
root
A
B
D
E
level 1
C
level 2
F
level 3
G
level 4
Fig. 4. An example of a tree.
generally called a push operation. Removing an element from a stack is called a pop operation. A queue is a sequence of components that allows component access and removal at one end (called the front of a queue) and addition of new components at the other end of the sequence (called the rear of a queue) [Fig. 3]. A tree is a data structure whose components are arranged in a hierarchy of levels. The topmost level has only one component called the root of the tree, but other levels can have many components. Each component on some level (except the root) is a child of some component on the level above. For example, each component on the second level is a child of the root, and each component on the third level is a child of some component on the second level, and so on. In summary, every component of a tree has one parent located on the level above (except the root, which does not have a parent) and possibly several children located on the level below. The number of levels in such a hierarchy is called the height of a tree. For example, the components A, B, C, D, E, F, and G will form a tree of height 4 if they are related in the following manner: A is the parent of B and C, B is the parent of D and E, C is the parent of F, and F is the parent of G (Fig. 4). A binary tree is a tree in which each component may have up to two children.
Fig. 5. A graph with four vertices and five edges.
A graph is similar to a tree in that it represents a collection of components and their relationship. However, the relationship among its components is not necessarily hierarchical, and any component can be related to any number of components. A component is called a vertex or node, and a relation between two vertices is called an edge. A convenient way to show a graph is to draw a picture in which a vertex is represented by a dot and an edge is represented by a line connecting two vertices (Fig. 5). If the edges have a direction assigned to them, the graph is called a directed graph. See GRAPH THEORY. Elementary algorithms. Many algorithms are useful in a broad spectrum of computer applications. These elementary algorithms are widely studied and considered an essential component of computer science. They include algorithms for sorting, searching, text processing, solving graph problems, solving basic geometric problems, displaying graphics, and performing common mathematical calculations. Sorting arranges data objects in a specific order, for example, in numerically ascending or descending orders. Internal sorting arranges data stored internally in the memory of a computer. Simple algorithms for sorting by selection, by exchange, or by insertion are easy to understand and straightforward to code. However, when the number of objects to be sorted is large, the simple algorithms are usually too slow, and a more sophisticated algorithm, such as heap sort or quick sort, can be used to attain acceptable performance. External sorting arranges data records stored on disk or tape. Searching looks for a desired data object in a collection of data objects. Elementary searching algorithms include linear search and binary search. Linear search examines a sequence of data objects one by one. Binary search adopts a more sophisticated strategy and is faster than linear search when searching a large array. A collection of data objects that are to be frequently searched can also be stored as a tree. If such a tree is appropriately structured, searching the tree will be quite efficient. A text string is a sequence of characters. Efficient algorithms for manipulating text strings, such as algorithms to organize text data into lines and paragraphs and to search for occurrences of a given pattern in a document, are essential in a word processing system. A source program in a high-level programming language is a text string, and text processing is a necessary task of a compiler. A compiler needs to use efficient algorithms for lexical analysis (grouping individual characters into meaningful words or symbols) and parsing (recognizing the syntactical structure of a source program). See SOFTWARE ENGINEERING; WORD PROCESSING. A graph is useful for modeling a group of interconnected objects, such as a set of locations connected by routes for transportation. Graph algorithms are useful for solving those problems that deal with objects and their connections—for example, determining whether all of the locations are connected, visiting all of the locations that can be reached from a
447
448
Algorithm given location, or finding the shortest path from one location to another. Geometric objects such as line segments, circles, and polygons are fundamental in modeling many physical objects, such as a building or an automobile. Geometric algorithms are useful for solving many problems in designing and analyzing a geometric model of a physical object. Basic geometric algorithms include those for determining whether two line segments intersect, and whether a given point is within a given polygon. A related area, computer graphics, is the study of methods to represent and manipulate images, and to draw them on a computer screen or other output devices. Basic algorithms in the area include those for drawing line segments, polygons, and ovals. See COMPUTER GRAPHICS. Mathematical algorithms are of wide application in science and engineering. Basic algorithms for mathematical computation include those for generating random numbers, performing operations on matrices, solving simultaneous equations, and numerical integration. Modern programming languages usually provide predefined functions for many common computations, such as random number generation, logarithm, exponentiation, and trigonometric functions. Recursion. An algorithm for solving a problem can often be naturally expressed in terms of using the algorithm itself to solve a simpler instance of the problem. Euclid’s algorithm makes use of the fact that the greatest common divisor of two different integers is the same as the greatest common divisor of the smaller number of the two and their difference. By recognizing that the problem of finding the greatest common divisor of the smaller number and the difference is a simpler instance of the given problem, the algorithm can be stated recursively: If the two given numbers are equal, their value is the answer; otherwise, (this is the recursive part) find the greatest common divisor of the smaller number and their difference. The following shows the steps of this recursive algorithm when used to find the greatest common divisor of 48 and 18. a. Since 18 is smaller, find the greatest common divisor of 18 and 30 (= 48 − 18). b. Since 18 is smaller, find the greatest common divisor of 18 and 12 (= 30 − 18). c. Since 12 is smaller, find the greatest common divisor of 12 and 6 (= 18 − 12). d. Since 6 is smaller, find the greatest common divisor of 6 and 6 (= 12 − 6). e. The greatest common divisor is 6. Since 6 is the greatest common divisor of 6 and 6, it is also the greatest common divisor of 6 and 12, of 18 and 12, of 18 and 30, and of the given numbers 48 and 18. Analysis of algorithms. A problem to be solved by an algorithm generally has a natural problem size, which is usually the amount or the magnitude of
data to be processed, and the number of operations that the algorithm performs depends predominantly on the problem size. A common approach to characterize the time efficiency of an algorithm is to count the number of operations that the algorithm performs in the worst case. For example, the problem of searching an array A1, A2, . . . , An for a value V has the problem size n. The linear search algorithm successively compares the desired value V with the array components A1, A2, . . . , An until it finds that a component is equal to V or until it has compared V with all of the components without finding one that is equal to V. The worst case obviously requires the algorithm to compare V with all n components. Therefore, the number of operations performed by linear search is essentially proportional to n, and linear search is said to have a running time of O(n). If the array is sorted, binary search can be used instead of linear search. Binary search first compares V with the middle component Am, where m is (1 + n)/2. If Am is equal to V, it terminates; otherwise, the search continues with at most one-half of the array components remaining to be considered. (Depending on whether V is greater than or less than Am, either the components A1 . . . Am−1 or the components Am+1 . . . An still need to be considered in the search, since the array is sorted.) A comparison of V with the middle component of the part of the array remaining to be considered either will lead to success of the search and termination of the algorithm, or again will reduce the number of elements to be searched by at least one-half of those remaining to be considered, and so forth. In summary, each comparison either leads to termination of the algorithm or reduces the number of components remaining to be considered in the search by at least one-half. In the worst case, about log2n comparisons are needed to reduce the number of components remaining in the search to one (by then the search is trivial). The number of operations is essentially proportional to log2n, and the algorithm is said to have a running time of O(log n). When the problem size n is large, log2n is much smaller than n, indicating that binary search will be much faster than linear search when searching a large array. An open question in the study of computational complexity concerns a class known as NP-complete problems. The class includes such problems as the traveling salesperson problem (minimizing the distance traveled in a tour of a number of cities) and finding a hamiltonian path in a graph (a path that passes through each vertex once). Although the correctness of proposed solution to these problems can be checked in polynomial time (a running time that is a polynomial function of the problem size), no polynomial-time algorithms to solve these problems have been found. (Known algorithms to solve them have an exponential running time, which makes them impractical when the problem size is large.) Furthermore, if any NP-complete problem can be solved in polynomial time, so can all others in the
Alicyclic hydrocarbon class. Whether these problems can be solved in polynomial time remains an open question. See CRYPTOGRAPHY. Randomized algorithms. By using a randomizer (such as a random number generator) in making some decisions in an algorithm, it is possible to design randomized algorithms that are relatively simple and fast. Since the output of a randomizer may vary from run to run, the output and the running time of a randomized algorithm may differ from run to run for the same input. There is a nonzero probability that the output of a randomized algorithm may be incorrect and execution of the algorithm may be slower than desired. A design goal of randomized algorithms is to keep such a probability small. However, there are critical systems that cannot tolerate a nonzero, albeit small, probability of incorrect output or a longer running time. Adaptive systems. In many applications, a computer program needs to adapt to changes in its environment and continue to perform well. An approach to make a computer program adaptive is to use a self-organizing data structure, such as one that is reorganized regularly so that those components most likely to be accessed are placed where they can be most efficiently accessed. A self-modifying algorithm that adapts itself is also conceivable. For developing adaptive computer programs, biological evolution has been a source of ideas and has inspired evolutionary computation methods such as genetic algorithms. See GENETIC ALGORITHMS. Parallel algorithms. Certain applications require a tremendous amount of computation to be performed in a timely fashion. An approach to save time is to develop a parallel algorithm that solves a given problem by using a number of processors simultaneously. The basic idea is to divide the given problem into subproblems and use each processor to solve a subproblem. The processors usually need to communicate among themselves so that they may cooperate. The processors may share memory, through which they can communicate, or they may be connected by communication links into some type of network such as a hypercube. See CONCURRENT PROCESSING; MULTIPROCESSING; SUPERCOMPUTER. Samuel C. Hsieh Bibliography. F. M. Carrano, P. Helman, and R. Veroff, Data Abstraction and Problem Solving with C++, Addison Wesley Longman, 1998; E. Horowitz, S. Sahni, and S. Rajasekaran, Computer Algorithms/ C++, Computer Science Press, 1996; B. R. Preiss, Data Structures and Algorithms with ObjectOriented Design Patterns in C++, Wiley, 1999.
rings. Compounds with one to five alicyclic rings of great variety and complexity are found in many natural products such as steroids and terpenes. By far the majority of these have six-membered rings. See AROMATIC HYDROCARBON; STEROID; TERPENE. Structures and nomenclature. The bonding in cyclic hydrocarbons is much the same as that in open-chain alkanes and alkenes. An important difference, however, is the fact that the atoms in a ring are part of a closed loop. Complete freedom of rotation about a carbon-carbon bond (C C) is not possible; the ring has faces or sides, like those of a plate. Simple monocyclic hydrocarbons are usually represented as bond line structures; for example, cyclopropane (1a) is usually represented as structure (1b) and methylcycloheptane as structure (2). These CH2 H2C
CH2
(1a)
(1b) CH3
(2) hydrocarbons are named by adding the prefix cyclo to the stem of the alkane corresponding to the number of atoms in the ring. When two or more substituents are attached to the ring, the relative positions and orientation must be specified: cis on the same side and trans on the other, as in cis-1,2-dimethylcyclopentane (3) and trans-1,3dimethylcyclopentane (4). CH3
CH3 CH3
CH3
(3)
(4)
In bicyclic compounds the rings can be joined in three ways: spirocyclic, fused, and bridged, as illustrated in the structures for spiro[4.5]decane (5), bicyclo[4.3.0]nonane (6), and bicyclo[2.2.1]heptane (7). In each case the name indicates the total 2
1
6
7
5
5 3
4
10
4
6
3
1
8
9
2
(5)
(6) 7
Alicyclic hydrocarbon An organic compound that contains one or more closed rings of carbon atoms. The term alicyclic specifically excludes carbocyclic compounds with an array of π -electrons characteristic of aromatic
7
8
1
6 5
2 4
(7)
3
9
449
450
Alicyclic hydrocarbon number of carbon atoms,and the number of atoms in each bridge. Atoms are numbered as shown. Any of these bicyclic systems can be transformed to a tricyclic array by introduction of another bond between nonadjacent carbons, as in structure (8), or an additional ring, as in structure (9). Cyclic structure
H
(a)
eclipsed
H
(b)
Fig. 2. Nonplanar (a) chair and (b) boat structures of cyclohexane.
(8)
(9) gives rise to the possibility of compounds made up of molecular subunits that are linked mechanically rather than chemically. In rotaxanes (10), bulky groups are introduced at the ends of a long chain that is threaded through a large ring (>C30). Cyclization of the ends leads to a catenane (11). Several
(10)
tor is the interaction that occurs when bonds on adjacent carbons are aligned (eclipsing interaction; Fig. 2), as in cyclopropane or the boat form of cyclohexane. The rings in cyclobutane and cyclopentane are bent slightly to reduce these eclipsing interactions; cyclohexane exists almost entirely in a somewhat twisted chair conformation. In medium rings of 8–11 carbon atoms, another destabilizing effect is the interaction of atoms situated across the ring from each other. The presence of a cis double bond in small rings increases the ring strain somewhat, but cyclobutene and even cyclopropene can be prepared. A cyclopropene ring is present in the naturally occurring sterculic acid (12 ). Alicyclic rings with a trans double bond or a triple bond are possible when the ring contains eight or more carbons, as in transcyclooctene (13).
(11)
examples of compounds with these structures have been prepared.
CH3(CH2)7
(CH2)7CO2H
(12)
Ring strain and conformation. In a planar ring of n carbon atoms, the internal bond angles are given by n−2 180 n and the angles (α) for three-, four-, and fivemembered rings are 60◦, 90◦, and 108◦, respectively. The angle between bonds of a tetrahedral carbon atom is 109◦. Thus in the smaller rings, particularly cyclopropane, a significant distortion of the tetrahedral bond angle is required, and the ring is strained. (Fig. 1). See BOND ANGLE AND DISTANCE.
60° (a)
90° (b)
109°
108° (c)
(d)
Fig. 1. Internal bond angles (α) for (a–c) planar rings and (d) a tetrahedral carbon atom.
Six-membered and larger rings can have normal bond angles by adopting nonplanar puckered conformations. For cyclohexane, two nonplanar arrangements are called chair and boat (Fig. 2). See CONFORMATIONAL ANALYSIS. Angle strain is not the only consideration in the conformation of alicyclic compounds. Another fac-
(13) Properties. The boiling points, melting points, and densities of cycloalkanes are all higher than those of their open-chain counterparts, reflecting the more compact structures and greater association in both liquid and solid. Geometrical constraints in the smaller rings have significant effects on reactions of cycloalkanes and derivatives. Because of ring strain, ring-opening reactions of cyclopropane, such as isomerization to propene, take place under conditions that do not affect alkanes or larger-ring cycloalkanes [reaction (1)]. Comparison of reaction rates of CH3HC
CH2
(1)
compounds with different ring size or cis-trans configuration has provided important insights about reaction mechanism and conformational analysis. See ORGANIC REACTION MECHANISM. Ring-forming reactions. A number of useful methods have been devised that lead to alicyclic rings. These reactions are of three types: C-C bond formation between atoms in an open-chain precursor,
Alicyclic hydrocarbon cycloaddition or cyclooligomerization, and expansion or contraction of a more readily available ring. In cycloaddition, two molecules react with formation of two bonds; in cyclooligomerization, three or more molecules combine to form three or more bonds. Three-membered rings. Cyclopropane can be prepared by dehalogenation of 1-bromo-3-chloropropane (14) with zinc (Zn), as in reaction (2). A more general BrCH2CH2CH2Cl
OH CN
O HCN
. HNO 2
lH 4, 2
1. LiA
O–
(6) O
+
CH2N2
–N2
Zn
(2)
(14) approach is cycloaddition of a carbene or carbenoid reagent to a double bond, as in reactions (3).
Eight-membered rings. The compound of central importance in this series is cyclooctatetraene (15), which is derived from four molecules of acetylene by cyclooligomerization with a nickel (Ni) catalyst [reaction (7)]. Similarly, dimerization of 1,3-dienes
4 HC C
C
CCl2
Cl
CH Ni
CH2I2 - Zn/Cu
Cl
CHCl
H
(7)
(3)
Cl
(15) Four-membered rings. One of the important routes to
cyclobutanes is cycloaddition of two double bonds. These reactions are facile with allenes, fluoroalkenes, and other alkenes in which a free-radical intermediate can be stabilized, as in reaction (4). 2+2 Cycload-
gives cycloocta-1,5-dienes [reaction (8)] and trimerization leads to cyclododecatrienes [reaction (9)]. H2C
CH
CH
CH2
(C6H5)
(8) C6H5HC
CH2 + F2C
CCl2
Cl
F Cl
F
ditions (in which a two-atom molecule combines with a two-atom partner to give a four-atom ring) also occur photochemically [reaction (5)]. FourO h
+ H2C
C(CH3)2
(9)
(4)
Cyclooctatetraene is a highly reactive, nonaromatic polyene that exists in a tub conformation. In several reactions, the products are those arising from the bicyclic valence tautomer (an isomer formed rapidly and reversibly by changes in valence bonds), as in reaction (10).
(5)
(10)
CH3 CH3
membered rings can also be obtained by ring contraction reactions of cyclopentanes. Five- and six-membered rings. Compounds with these rings are available in quantity, and there are many ways to interconvert them. Intramolecular reactions between groups separated by four or five atoms are highly favorable; an example is the cationic cyclization in the biosynthesis of terpenes and steroids. Another very general route to cyclohexane derivatives is Diels-Alder cycloaddition. See BIOSYNTHESIS; DIELSALDER REACTION. Seven-membered rings. Cycloheptane derivatives can be obtained by ring expansion reactions such as those in reactions (6).
See TAUTOMERISM. Larger rings. Ring closure between ester or nitrile groups at the ends of a chain by base-catalyzed condensation can be used to obtain rings of 14 atoms or more, as in reaction (11), where Et = C2H5. Very :B H CO2Et
CO2Et
(11)
O C
O OEt
dilute solutions must be used in order to minimize
451
452
Alismatales intermolecular reaction. The yields of medium rings (9–12 ring atoms) by these reactions are negligible, however, because the geometry necessary for such an enolate addition to the ester is unfavorable. A useful method for all rings is known as the acyloin condensation of a diester on molten sodium (Na), as in reaction (12), where R = an organic group. In this CO2R CO2R
ONa Na
(12)
Alismatidae
ONa
process, the ends of the chain can be aligned by moving on the sodium surface. Major compounds. Two alicyclic compounds are manufactured in large volume; both are products of the petroleum industry. Cyclopentadiene (16) is formed from various alkylcyclopentanes in naphtha fractions during the refining process. It is a highly reactive diene and spontaneously dimerizes in a 4 + 2 cycloaddition [reaction (13)]; it is used as a copolymer in several resins.
+
to form a near-basal side branch which has undergone its own sort of specialization while remaining relatively primitive in gross floral structure. The Alismatales and some related orders have often been treated as a single order Helobiae or Helobiales, embracing most of what is here treated as the subclass Alismatidae. See ALISMATIDAE; LILIOPSIDA; MAGNOLIOPHYTA; PLANT KINGDOM. Arthur Cronquist
(13)
(16) Cyclohexane is produced in large quantity by hydrogenation of benzene. The principal use of cyclohexane is conversion by oxidation in air to a mixture of cyclohexanol and the ketone, which is then oxidized further to adipic acid for the manufacture of nylon. See ORGANIC SYNTHESIS. James A. Moore Bibliography. F. A. Carey and R. J. Sundberg, Advanced Organic Chemistry: Structure and Mechanisms (Part A), 4th ed., 2004; M. Grossel, Alicyclic Chemistry, 1997; Z. Rappaport and S. Patai (eds.), Chemistry of Alkanes and Cycloalkanes, 1992.
Alismatales A small order of flowering plants, division Magnoliophyta (Angiospermae), which gives its name to the subclass Alismatidae of the class Liliopsida (monocotyledons). It consists of three families (Alismataceae, Butomaceae, and Limnocharitaceae) and less than a hundred species. They are aquatic and semiaquatic herbs with a well-developed, biseriate perianth that is usually differentiated into three sepals and three petals, and with a gynoecium of several or many, more or less separate carpels. Each flower is usually subtended by a bract. Butomus umbellatus (flowering rush) and species of Sagittaria (arrowhead, family Alismataceae) of this order are sometimes cultivated as ornamentals. The Alismatales are usually regarded as the most primitive existing group of monocotyledons, but because of their trinucleate pollen and nonendospermous seeds they cannot be considered as directly ancestral to the rest of the group. Instead they appear
A relatively primitive subclass of the class Liliopsida (monocotyledons) of the division Magnoliophyta (Angiospermae), the flowering plants, consisting of 4 orders, 16 families, and less than 500 species. Typically they are aquatic or semiaquatic, with apocarpous flowers and nonendospermous seeds. They have trinucleate pollen, and the stomates usually have two subsidiary cells. The orders Alismatales, Hydrocharitales, and Najadales are closely related among themselves and have often been treated as a single order, Helobiae or Helobiales. The Triuridales differ from the other orders in being terrestrial and mycotrophic, without chlorophyll, and in having abundant endosperm in the seeds. See ALISMATALES; HYDROCHARITALES; TRIURIDALES; MAGNOLIOPHYTA; NAJADALES; PLANT KINGDOM. Arthur Cronquist; T. M. Barkley
Alkali Any compound having highly basic properties, strong acrid taste, and ability to neutralize acids. Aqueous solutions of alkalies are high in hydroxyl ions, have a pH above 7, and turn litmus paper from red to blue. Caustic alkalies include sodium hydroxide (caustic soda), the sixth-largest-volume chemical produced in the United States, and potassium hydroxide. They are extremely destructive to human tissue; external burns should be washed with large amounts of water. The milder alkalies are the carbonates of the alkali metals; these include the industrially important sodium carbonate (soda ash) and potassium carbonate (potash), as well as the carbonates of lithium, rubidium, and cesium, and the volatile ammonium hydroxide. Sodium bicarbonate is a still milder alkaline material. See ACID AND BASE; PH. Soda ash and potash washed from the ashes of wood or other biomass fires were among the first manufactured chemicals, and they produced crude soap when reacted with fats. Today, the manufacture of soap and synthetic detergents still requires large amounts of alkalies. See DETERGENT; SOAP. Synthetic soda ash had been made from salt since about 1840, with the Solvay process, using limestone and ammonia, predominating since 1950. Beginning about 1940, large sources of mineral deposits of impure soda ash crystallized from ancient lakes were mined and purified in increasing quantities; however, in the 1980s, production of synthetic and natural soda ash became about equal.
Alkali emissions
Alkali emissions Light emissions in the upper atmosphere from elemental lithium, potassium, and especially sodium. These alkali metals are present in the upper atmosphere at altitudes from about 50 to 62 mi (80 to 100 km) and are very efficient in resonant scattering of sunlight. The vertical column contents (number of atoms per square meter) of the alkali atoms are easily deduced from their respective emission intensities. First detected with ground-based spectrographs, the emissions were observed mainly at twilight since they tend to be overwhelmed by intense scattered sunlight present in the daytime. A chemiluminescent process gives rise to so-called nightglow emissions at the same wavelengths. The development of lidars (laser radars) that are tuned to the resonance lines have enabled accurate resolution of the concentrations of these elements versus altitude for any time of the day. See AERONOMY; CHEMILUMINESCENCE. There is little doubt that the origin of these metals is meteoritic ablation. Rocket-borne mass spectrometers have found that meteoritic ions are prevalent above the peaks of neutral atoms with a composition similar to that found in carbonaceous chondrites, a common form of meteorites. The ratio of the concentrations of ions to neutral atoms rises rapidly with altitude above 55 mi (90 km). See METEORITE. Sodium emission. Sodium (Na) is the most abundant alkali metal in meteorites. The sodium D-lines, a doublet at 589 and 589.6 nanometers, were first detected in the nightglow in the late 1920s and at twilight a decade later. The nominal peak concen-
100
60
95
altitude, mi
altitude, km
Until about 1900, all caustic soda was made chemically from soda ash by action with slaked lime, and this method was used for most of the soda ash manufactured until about World War II. Meanwhile, the availability of inexpensive electric power had allowed the commercial development of salt electrolysis. This process yields (at the cathode) a 10– 12% sodium hydroxide solution and hydrogen gas and chlorine gas (at the anode). Since demand for chlorine increased rapidly after 1940, large quantities of its chemical equivalent, caustic soda, were being produced and sold. Subsequent industrial problems caused wide market swings involving the replacement of soda ash by caustic soda, and the marketing and pricing of soda ash, caustic soda, and chlorine. About 50% of the caustic soda produced goes into making many chemical products, about 16% into pulp and paper, 6.5% each into aluminum, petroleum, and textiles (including rayon), with smaller percentages into soap and synthetic detergents, and cellophane. For soda ash, about 50% goes to react mainly with sand in making glass, 25% to making miscellaneous chemicals, 6.5% each to alkaline cleaners and pulp and paper, and a few percent to water treatment and other uses. See ALKALI METALS; ELECTROCHEMICAL PROCESS; GLASS; HYDROXIDE; PAPER; SOAP; TEXTILE CHEMISTRY. Donald F. Othmer
Na
Na+
55
90
85 0
1000
2000
50
concentration, cm −3 Fig. 1. Summer observations of sodium (Na) and sodium ion (Na+) at high latitudes. ( After W. Swinder, Sodium chemistry: A brief review and two new mechanisms for sudden sodium layers, Planet. Space Sci., 40: 247–253, 1992 )
tration of sodium is 3 × 109 atoms m−3 near 55 mi (90 km) where the total gas concentration of the atmosphere is 7 × 1019 atoms (or molecules) m−3. The sodium concentration declines by half at 3 mi (5 km) above and below the peak. The altitude of the maximum concentration for the sodium layer is in accord with the fact that the bulk of the micrometeorites ablate at this height, since their mean velocity is about 22 mi/s (35 km/s). Mesospheric dynamics commonly distorts the sodium profile. Lidar measurements have provided the best evidence of these distortions and will allow scientists to achieve a better understanding of the dynamics of the mesosphere. See MICROMETEORITE. The discovery of the occasional appearance of socalled sudden sodium layers (Fig. 1) is indicative of the complex dynamics and chemistry in the upper mesosphere. While some fresh meteor trails may conceivably contribute to such narrow layers of neutral meteor atoms, it is likely that the principal mechanism for sudden sodium layers originates with the neutralization of meteoritic ions, which can be layered by a process involving the winds and the Earth’s magnetic field. Lidar has established that the vertical column content of sodium, about 3 × 1013 atoms m−2 at midlatitudes during equinox, varies very little diurnally. Twilight spectroscopic data, confirmed by the more recent lidar data, have long indicated that there is a seasonal variation in sodium, with more present in winter. The variation increases with latitude, about a factor of 2 near 20◦ latitude and a factor of 5 near 55◦ (Fig. 2). Data for polar latitudes are scarce. This pattern is thought to relate to the chemistry. Thus, the main depletion process for sodium, reaction (1), Na + O2 + M → NaO2 + M
(1)
where M is either oxygen (O2) or nitrogen (N2), proceeds more slowly as temperature increases. The temperature at mesospheric altitudes is greater in winter than in summer. See LIDAR. Other factors may be more or equally important, like the seasonal variation of atomic oxygen. The decline of sodium below its peak would be even
453
Alkali metals
8
8
6
6 sodium 4
4
2
2 potassium
0
June July Aug Sept Oct Nov Dec Jan Feb Mar Apr MayJune July Aug Sept
vertical column content (K), 1012 m −2
vertical column content (Na), 1013 m −2
454
0
The lithium nightglow emission is undetectable. See ALKALI METALS; ATMOSPHERIC CHEMISTRY; MESOSPHERE. William Swider Bibliography. P. P. Batista et al., Characteristics of the sporadic sodium layers observed at 23◦S, J. Geophys. Res., 94:15349–15358, 1989; J. M. C. Plane, The chemistry of meteoric metals in the Earth’s upper atmosphere, Int. Rev. Phys. Chem., 10:55–106, 1991; W. Swider, Sodium chemistry: A brief review and two new mechanisms for sudden sodium layers, Planet. Space Sci., 40:247–253, 1992.
Alkali metals
month (1975 -1976) ◦
Fig. 2. Seasonal patterns of the vertical column contents of sodium and potassium at 44 latitude. (After W. Swider, Chemistry of mesospheric potassium and its different seasonal behavior as compared to sodium, J. Geophys. Res., 92(D5):5621–5626, 1987)
sharper if it were not for reactions (2a) and (2b). NaO2 + O → NaO + O2
(2a)
NaO + O → Na + O2
(2b)
Process (2b) can leave sodium in an excited˜state. Relaxation of the excited atoms to the ground state produces the sodium D-lines. The father of aeronomy, S. Chapman, proposed this nightglow mechanism a half century ago. Current models of the nightglow yield emission intensities compatible with observations. Potassium emission. Potassium (K) is 15 times less abundant than sodium in meteorites. The ratio of the potassium and sodium column contents in the mesosphere ranges from 1/10 to 1/100 because, unlike sodium, the column content of potassium, 2 × 1012 m−2, is fairly constant with season (Fig. 2). Once thought related to a possible additional source for sodium like sea spray, this variation may have a chemical origin. Although specific reactions may be equivalent, the particular rate coefficients may differ. For example, the rate coefficient for the formation of potassium superoxide (KO2) is about twice that for sodium superoxide (NaO2), and is less strongly dependent on temperature. The potassium doublet at 767 and 770 nm is estimated to have a nightglow intensity near the night sky background, 50 times smaller than the typical intensity of the sodium nightglow. The potassium nightglow has never been detected. Lithium emission. Lithium (Li) is 35 times less abundant in meteorites than potassium, and 500 times less abundant than sodium. The meteoric source of lithium can be swamped at times by other sources, such as nuclear explosions, volcanic eruptions, and deliberate releases of lithium for aeronomic investigations of the upper atmosphere. Nevertheless, the lithium emission at 671 nm has been observed at twilight by spectrometers and at night by lidars, which indicate that the vertical content, 1010 m−2 at equinox, increases in winter like sodium.
The elements of group 1 in the periodic table. Of the alkali metals, lithium differs most from the rest of the group, and tends to resemble the alkaline-earth metals (group 2 of the periodic table) in many ways. In this respect lithium behaves as do many other elements that are the first members of groups in the periodic table; these tend to resemble the elements in the group to the right rather than those in the same group. Francium, the heaviest of the alkali-metal
Isotopes of the alkali metals
Element
Normal at. wt
Lithium, Li
6.939
Sodium, Na
22.9898
Potassium, K
39.102
Rubidium, Rb 85.47
Cesium, Cs
Francium, Fr
Mass no.
Radioactive
5 6 7 8 9 20 21 22 23 24 25 37 38 39 40 41 42 43 44 81 82 83 84 85 86 87
Yes No No Yes Yes Yes Yes Yes No Yes Yes Yes Yes No Yes No Yes Yes Yes Yes Yes Yes Yes No Yes Yes
88 89 90 91 132.905 127–132 133 134 135 136 137 138–145 223
Yes Yes Yes Yes Yes No Yes Yes Yes Yes Yes Yes
Half-life∗ 10⫺21 s Stable (7.5) Stable (92.5) 0.83 s 0.17 s 0.23 s 23.0 s 2.6 y Stable (100) 15.0 h 60 s 1.2 s 7.7 m Stable (93.1) 1.2 ⫻ 109 y Stable (6.9) 12.4 h 22 h 27 m 4.7 h 6.3 h 80 d 23 m Stable (72.2) 19 d 6.2 ⫻ 1010 y (27.8) 18 m 15 m 2.7 m 14 m Short Stable (100) 3 ⫻ 106 y 2.3 y 1.3 d 37 y Short 21 m
∗Figures in parentheses indicate the percentage occurrence in nature.
Alkaloid elements, has no stable isotopes and exists only in radioactive form (see table). See PERIODIC TABLE. In general, the alkali metals are soft, low-melting, reactive metals. This reactivity accounts for the fact that they are never found uncombined in nature but are always in chemical combination with other elements. This reactivity also accounts for the fact that they have no utility as structural metals (with the possible exception of lithium in alloys) and that they are used as chemical reactants in industry rather than as metals in the usual sense. The reactivity in the alkali-metal series increases in general with increase in atomic weight from lithium to cesium. See CESIUM; ELECTROCHEMICAL SERIES; FRANCIUM; LITHIUM; POTASSIUM; RUBIDIUM; SODIUM. Marshall Sittig Bibliography. H. V. Borgstedt and C. K. Matthews, Applied Chemistry of the Alkali Metals, 1987; J. Bowser, Inorganic Chemistry, 1993; A. S. Kertes and C. A. Vincent (eds.), Alkali Metal, AlkalineEarth Metal and Ammonium Halides in Amide Solvents, 1980.
cium. All the metals dissolve readily in acid. Whereas calcium reacts smoothly with water to yield hydrogen, the heavier members react as violently as sodium does. All the metals are soluble in liquid ammonia, yielding strongly reducing, electrically conducting, blue solutions. The order of solubility in water for most salts is calcium > strontium > barium, except that the order is reversed for the fluorides, hydroxides, and oxalates. All three elements unite directly with hydrogen to form hydrides and with nitrogen to form nitrides, but whereas ease of formation of a nitride increases with atomic number, ease of formation of a hydride decreases. The elements and their compounds find important industrial uses in low-melting alloys, deoxidizers, and drying agents and as cheap sources of alkalinity. See BARIUM; BERYLLIUM; CALCIUM; MAGNESIUM; PERIODIC TABLE; RADIUM; STRONTIUM. Reed F. Riley Bibliography. F. A. Cotton et al., Advanced Inorganic Chemistry, 6th ed., Wiley-Interscience, 1999; T. R. Dickson, Introduction to Chemistry, 7th ed., 1995; K. Mackay, Introduction to Modern Inorganic Chemistry, 6th ed., 2002.
Alkaline-earth metals Usually calcium, strontium, and barium, the heaviest members of group 2 of the periodic table (excepting radium). Other members of the group are beryllium, magnesium, and radium, sometimes included among the alkaline-earth metals. Beryllium resembles aluminum more than any other element, and magnesium behaves more like zinc and cadmium. The gap between beryllium and magnesium and the remainder of the elements of group 2 makes it desirable to discuss these elements separately. Radium is often treated separately because of its radioactivity. J. J. Berzelius first reduced the three alkaline-earth metals to the elementary state, but obtained them as amalgams by electrolysis. Humphry Davy in 1808 isolated the metals in the pure state by distillation of amalgams produced electrolytically. Today industrial preparation of these elements involves electrolysis of their molten chlorides or reduction of their oxides with aluminum. The alkaline earths form a closely related group of highly metallic elements in which there is a regular gradation of properties. The metals, none of which occurs free in nature, are all harder than potassium or sodium, softer than magnesium or beryllium, and about as hard as lead. The metals are somewhat brittle, but are malleable, extrudable, and machinable. They conduct electricity well; the specific conductivity of calcium is 45% of that of silver. The oxidation potentials of the triad are as great as those of the alkali metals. The alkaline earths exist as large divalent cations in all their compounds, in which the elements are present in the 2+ oxidation state. The metals have a gray-white luster when cut but tarnish readily in air. They burn brilliantly in air when heated, and form the metal monoxide, except for barium, which forms the peroxide. A certain amount of nitride is formed simultaneously, especially with cal-
Alkaloid A cyclic organic compound that contains nitrogen in a negative oxidation state and is of limited distribution among living organisms. Over 10,000 alkaloids of many different structural types are known; and no other class of natural products possesses such an enormous variety of structures. Therefore, alkaloids are difficult to differentiate from other types of organic nitrogen-containing compounds. Simple low-molecular-weight derivatives of ammonia, as well as polyamines and acyclic amides, are not considered alkaloids because they lack a cyclic structure in some part of the molecule. Amines, amine oxides, amides, and quaternary ammonium salts are included in the alkaloid group because their nitrogen is in a negative oxidation state (the oxidation state designates the positive or negative character of atoms in a molecule). Nitro and nitroso compounds are excluded as alkaloids. The almost-ubiquitous nitrogenous compounds, such as amino acids, amino sugars, peptides, proteins, nucleic acids, nucleotides, prophyrins, and vitamins, are not alkaloids. However, compounds that are exceptions to the classical-type definition (that is, a compound containing nitrogen, usually a cyclic amine, and occurring as a secondary metabolite), such as neutral alkaloids (colchicine, piperine), the β-phenyl-ethylanines, and the purine bases (caffeine, theophylline, theobromine), are accepted as alkaloids. Nomenclature. The nomenclature of alkaloids has not been systematized. Most alkaloids bear the suffix -ine and their names are derived in various ways: from the generic name of the plant, for example hydrastine from Hydrastis canadensis; from the specific plant name, for example cocaine from Erythroxylum coca; from the common name of the drug yielding
455
456
Alkaloid them, for example, ergotamine from ergot; or from a particular physiological activity exhibited, for example, morphine from Morpheus, Greek god of dreams. Only the pelletierine group is named for a person, P. J. Pelletier, who discovered a number of alkaloids in the first third of the nineteenth century. Occurrence. Alkaloids are derived from both plants and animals, including marine organisms. Plant-derived alkaloids. Alkaloids often occur as salts of plant acids such as malic, meconic, and quinic acids. Some plant alkaloids are combined with sugars, for example, solanine in potato (Solanum tuberosum) and tomatine in tomato (Lycopersicum esculentum). Others occur as amides, for example, piperine from black pepper (Piper nigrum), or as esters, for example, cocaine from coca leaves (Erythroxylum coca). Still other alkaloids occur as quaternary salts or tertiary amine oxides. Although about 40% of all plant families contain at least one alkaloid-bearing species, only about 9% of the over 10,000 plant genera produce alkaloids. Among the angiosperms, alkaloids occur in abundance in certain dicotyledons and especially in the families Apocynaceae (dogbane, quebracho), Asteraceae (groundsel, ragwort), Berberidaceae (European barberry), Fabaceae (broom, gorse, lupine), Loganiaceae (Strychnos species), Menispermaceae (moonseed), Papaveraceae (poppies, chelidonium), Ranunculaceae (aconite, larkspur), Rubiaceae (cinchona bark, ipecac), Rutaceae (citrus, fagara), and Solanaceae (tobacco, deadly nightshade, tomato, potato, thorn apple). Rarely are they present in cryptogamia, gymnosperms, or monocotyledons. Among the monocotyledons, the Amaryllidaceae (amaryllis, narcissus) and Liliaceae (meadow saffron, veratrums) are families that bear alkaloids. Examples of well-known alkaloids are morphine and codeine from the opium poppy (Papaver somniferum), strychnine from Strychnos species, quinine (1) from Cinchona bark and various Cinchona
species, and coniine from poison hemlock (Conium maculatum). Other important alkaloids are colchicine from the autumn crocus (Colchicum autumnale), caffeine from coffee, tea, kola, and mat´e, and nicotine from the tobacco plant (Nicotiana tabacum) and Aztec tobacco (N. rustica). See CAFFEINE; COFFEE; COLA; COLCHICINE; NICOTINE ALKALOIDS. Animal-derived alkaloids. While most alkaloids have been isolated from plants, a large number have been isolated from animal sources. They occur in mammals, anurans (frogs, toads), salamanders, arthro-
pods (ants, millipedes, ladybugs, beetles, butterflies), marine organisms, mosses, fungi, and certain bacteria. Samandarine has been isolated from the European fire salamander (Salamandra maculosa). Very toxic steroidal alkaloids occur in certain frogs and toads; for example, the Columbian arrow-poison frog (Phyllobates aurotaenia) produces a highly lethal venom containing batrachotoxin. The Canadian beaver (Castor fiberi) produces castoramine, and the scent gland of the musk deer (Moschus moschiferus) produces muscopyridine. The Southern fire ant (Solenopsis invicta) secretes a powerful venom containing a series of 2,6-dialkylpiperidines. The European millipede (Glomeris marginata) when provoked discharges a defensive secretion containing glomerine, a quinazolone. When molested, the European ladybug secretes hemolymph at the joints to afford protection against predators; this hemolymph contains precoccinelline and its oxide, coccinelline. Certain species of butterflies elaborate pheromones derived from exogenous pyrrolizidine alkaloid precursors; the pyrrolizidine alkaloid occurs in the secretions of Lycorea ceres ceres and Danaus gilippus berenice, where it elicits olfactory-receptor responses in female antennae and serves as a chemical messenger that induces mating behavior. See PHEROMONE. Marine-derived alkaloids. Many alkaloids have been isolated from marine organisms, both plant and animal. Of great interest are the alkaloids occurring in certain toxic dinoflagellates, such as Gonyaux tamarensis in the North Atlantic and G. cantanella in the South Atlantic. The toxic principles include the alkaloids saxitoxin, gonyautoxin-II, and gonyautoxinIII. These paralytic alkaloids also occur in the Alaska butter clam (Saxidomus giganteus) and toxic mussels (Mytilus californianus). Tetrodotoxin, present in the ovaries and liver of the Japanese puffer fish (Spheroides rubripes and S. vermicularis), in the Taiwanese goby fish (Gobius cringer), and in the California newt (Taricha torosa), is one of the most lethal low-molecular-weight toxins known. Isolation and purification. Since few plant species produce only a single alkaloid, there is the problem of separating complex mixtures. The alkaloid mixture is usually extracted from the finely powdered plant material by percolation with some solvent, such as methanol or aqueous ethanol. Evaporation leaves a residue, which is then dissolved in an organic solvent such as chloroform, methylene chloride, or toluene, and filtered. The solution is extracted repeatedly first with dilute sulfuric or hydrochloric acid and, after basification, with an organic solvent. These extracts when evaporated yield different mixtures of alkaloids which may be further separated by column chromatography over alumina or silica gel, vacuum liquid chromatography, countercurrent distribution, centrifugally accelerated radial thin-layer chromatography, or preparative-layer chromatography. Eventually, it is possible to obtain compounds in a high state of purity. See CHEMICAL SEPARATION TECHNIQUES; CHROMATOGRAPHY. Biosynthesis. Of great interest to chemists is the mode of alkaloid synthesis in plants. The
Alkane biosynthesis of most alkaloids can be derived hypothetically, by using a few well-known chemical reactions, from such relatively simple precursors as phenylalanine, tyrosine, tryptophan, histidine, methionine, acetate units, terpene units, and amino acids such as anthranilic acid, lysine, and ornithine. A number of simple alkaloids have been synthesized from amino acid derivatives under physiological conditions: pelletierine is derived from the amino acid lysine, while cocaine and nicotine are derived from ornithine (2). Many alkaloids, such as
(2) morphine, berberine, and emetine, incorporate a tetrahydroisoquinoline unit derived from tyrosine; and others, such as strychnine, quinine, reserpine, and camptothecin, are derived from tryptophan. The terpenoid and steroidal alkaloids, such as the diterpenoid alkaloid aconitine and the steroidal alkaloid tomatidine from the tomato, are likely derived from mevalonic acid lactone. See BIOSYNTHESIS; STEROID; TERPENE. Function. The exact function of alkaloids in plants is not well understood, but they are often regarded as by-products of plant metabolism. They are sometimes considered to be reservoirs for protein synthesis; as protective agents to discourage attack by animals, insects, or microbes; as regulators of growth, metabolism, and reproduction; and as end products of detoxification of substances whose accumulation might be injurious to the plant. Medicinals. Many alkaloids exhibit marked pharmacological activity, and some find important uses in medicine. Atropine, the optically inactive form of hyoscyamine, is used widely in medicine as an antidote to cholinesterase inhibitors such as physostigmine and insecticides of the organophosphate type; it is also used in drying cough secretions. Morphine and codeine are narcotic analgesics, and codeine is also an antitussive agent, less toxic and less habit forming than morphine. Colchicine, from the corms and seeds of the autumn crocus, is used as a gout suppressant. Caffeine, which occurs in coffee, tea, cocoa, and cola, is a central nervous system stimulant; it is used as a cardiac and respiratory stimulant and as an antidote to barbiturate and morphine poisoning. Emetine, the key alkaloid of ipecac root (Cephaelis ipecacuanha), is used in the treatment of amebic dysentery and other protozoal infections. Epinephrine (3), or adrenaline, produced in most
(3) animal species by the adrenal medulla, is used as a bronchodilator and cardiac stimulant and to counter allergic reactions, anesthesia, and cardiac arrest.
Pilocarpine, from Pilocarpus jaborandi, in hydrochloride or nitrate form is cholinergic and is used topically in the treatment of glaucoma; it is also used orally or subcutaneously to stimulate secretion of saliva in individuals who undergo therapy with ganglionic blocking agents. Quinine, used in many areas of the world to treat malaria, is used in the United States to treat chloroquine-resistant falciparum malaria. Quinidine, a stereoisomer of quinine, is an antiarrhythmic agent used to control auricular fibrillation. Reserpine, from Rauwolfia serpentina, produces sedation, tranquilization, and depression of blood pressure; it is sometimes used to decrease high blood pressure and as a tranquilizer in cases of nervousness, hysteria, and stress. Taxol, from the Pacific yew tree (Taxus brevifolia), is used in the treatment of ovarian and breast cancer. Tubocurarine chloride (curare, from Strychnos and Chondodendron species) is employed as a skelet al muscle relaxant during surgery and to control convulsions resulting from strychnine poisoning and tetanus. Vinblastine and vincristine, from the Madagascar periwinkle plant (Catharanthus roseus), are potent anticancer drugs; the former is used to treat generalized Hodgkin’s disease and choriocarcinoma, while the latter is used to treat acute leukemia. See CHEMOTHERAPY AND OTHER ANTINEOPLASTIC DRUGS; EPINEPHRINE; MORPHINE ALKALOIDS; PHARMACOGNOSY; QUININE; STRYCHNINE ALKALOIDS. S. William Pelletier Bibliography. A. Brossi and G. A. Cordell (eds.), The Alkaloids: Chemistry and Physiology, vol. 42, 1992; G. A. Cordell, Introduction to Alkaloids: A Biogenetic Approach, 1981; D. R. Dalton, The Alkaloids: The Fundamental Chemistry—A Biogenetic Approach, 1979; K. Mothes, H. R. Shutte, and M. Luckner, Biochemistry of Alkaloids, 1985; S. W. Pelletier (ed.), Alkaloids: Chemical and Biological Perspectives, vols. 1–9, 1983–1995; S. W. Pelletier, The nature and definition of an alkaloid, Alkaloids: Chemical and Biological Perspectives, vol. 1, 1983; T. Robinson, The Biochemistry of Alkaloids, 2d ed., 1981; I. W. Southon and J. Buckingham (eds.), Dictionary of Alkaloids, 1989.
Alkane An organic compound with the general formula CnH2n+2. Alkanes are open-chain (aliphatic or noncyclic) hydrocarbons with no multiple bonds or functional groups. They consist of tetrahedral carbon atoms, up to 105 carbons or more in length. The C-C σ bonds are formed from sp3 orbitals, and there is free rotation around the bond axis. See MOLECULAR ORBITAL THEORY; STRUCTURAL CHEMISTRY. Alkanes provide the parent names for all other aliphatic compounds in systematic nomenclature. Alkanes are designated by the ending -ane appended to a stem denoting the chain length. The straightchain isomer is designated by the prefix n(normal); other isomers are named by specifying the size of the branch and its location (see table). The number of isomers increases enormously in larger molecules;
457
458
Alkane Nomenclature and properties of alkanes
Name
Structure
Boiling point, ˚ C ( ˚ F)
Heat of formation ( Hf˚ ), kJ
Methane Ethane Propane n-Butane 2-Methylpropane (isobutane) n-Pentane 2-Methylbutane (isopentane) 2,2-Dimethyl propane (neopentane) n-Hexane
CH4 CH3 CH3 CH3 CH2 CH3 CH3 (CH2 )2 CH3 (CH3 )2 CHCH3 CH3 (CH2 )3 CH3 (CH3 )2 CHCH2 CH3
⫺162 (⫺260) ⫺89 (⫺128) ⫺42 (⫺44) ⫺0.5 (33) ⫺11.7 (10.9) 36.1 (97.0) 29.9 (85.8)
⫺74.5 ⫺83.45 ⫺104.6 ⫺125.7 ⫺134.2 ⫺146.87 ⫺153.7
9.4 (49) 68.7 (156)
⫺167.9 ⫺167.0
(CH3 )4 C CH3 (CH2 )4 CH3
thus there are 75 isomers of C10H22 and over 4 billion for C30H62. Properties. Alkanes with four or fewer carbons are gases at atmospheric pressure. Higher n-alkanes are liquids or, above about 20 carbons, solids known as paraffin wax. Alkanes have densities lower than that of water and have very low water solubility; other properties depend on the degree of branching (see table). The heat of formation ( H◦f) is a measure of the energy content of a compound relative to the component elements in standard states. For example, comparison of heat of formation values for three alkane isomers with the molecular formula C5H12 indicates that the relative energy content of the three pentane isomers decreases with increased branching; that is, the branched isomer is thermodynamically most stable. See PARAFFIN. Sources. Alkanes are the major components of natural gas and petroleum, which are the only significant sources. Much smaller amounts of alkanes have been produced from coal at various times and locations, either indirectly by the Fischer-Tropsch process or by direct liquefaction. See COAL LIQUEFACTION; FISCHER-TROPSCH PROCESS; NATURAL GAS; PETROLEUM. Certain long-chain alkanes and also some alicyclic hydrocarbons are very widely distributed both in sedimentary soils and shales and among living organisms. Phytane and pristane are derived from phytol, and these alkanes are of significance as biological markers in tracing diagenesis. Both insects and plants secrete alkanes, which function as moisture barriers and protective coatings. Cuticular waxes on fruit and leaves contain predominantly the n-alkanes with 29, 31, and 33 carbons; the preponderance of odd-carbon chains reflects the origin by loss of carbon monoxide (CO) or carbon dioxide (CO2) from even-carbon precursors. See DIAGENESIS. Individual lower alkanes can be separated from the more volatile distillate fractions of petroleum, but beyond the C7–C8 range the alkanes obtained are mixtures of many isomers. Compounds of a specific structure can be prepared on a laboratory scale by chemical synthesis. Various C-C bond-forming steps such as coupling or condensation are carried out to build up the desired carbon skeleton. The final step is usually removal of a functional group by some type of reduction. Several insect pheromones, for example 15,19,23-trimethylheptatriacontane (C40H82)
from the female tsetse fly, have been synthesized for biological study. See PHEROMONE. Chemistry. Much of the chemistry of alkanes begins at the petroleum refinery, where several reactions are carried out to adjust the hydrocarbon composition of crude oil to that needed for a constantly changing set of applications. Major reactions are (1) isomerization of straight-chain alkanes to branched compounds; (2) cracking to produce smaller molecules; (3) alkylation, for example, combination of propylene (1) and butane (2) to give 2,3-dimethylpentane (3), as in reaction (1); and CH3 CH
(1)
CH2 + CH3 CH2 CH2 CH3
(2) CH3 CH(CH3 )CH(CH3 )CH(CH3 )CH2 CH3
(1)
(3)
(4) cyclodehydrogenation (platforming), in which aromatizations occur. An important objective in some of these processes is to increase the yield of highly branched alkanes in the C6–C8 range needed for gasoline. See ALKYLATION (PETROLEUM); AROMATIZATION; CRACKING; GASOLINE. By far the most important end use of alkanes is combustion as fuel to provide heat and electric or motive power. In most cases, complete oxidation is not achieved, and varying amounts of incompletely oxidized fragments, carbon monoxide, and elemental carbon are produced. Controlled partial oxidation is possible if all the C-H bonds in an alkane are equivalent or if one C-H bond is significantly weaker than all the others. An example of the latter situation is isobutane, which is converted to the hydroperoxide on industrial scale for the manufacture of t-butyl alcohol. See AUTOXIDATION; COMBUSTION. Alkanes have been referred to as paraffin hydrocarbons to indicate their low affinity or reactivity. They contain no unshared electron pairs or accessible empty bonding orbitals, and they are unaffected by many reagents that attack π bonds or other functional groups. One type of reaction that does occur is substitution by a radical chain process. Examples are chlorination and vapor-phase nitration, involving the odd-electron species Cl· and NO2·, respectively. Neither reaction is selective; when two or more types of C-H bonds are present in the alkane, mixtures of products are usually obtained. Thus
Alkene propane gives rise to 1- and 2-chloropropanes as well as dichloro compounds. Nitration of propane leads to a mixture of 1- and 2-nitropropane, and also nitromethane and nitroethane by C-C bond cleavage. See HALOGENATED HYDROCARBON; HALOGENATION; NITRATION. Since alkanes are the most abundant class of organic compounds available, much effort has been directed to the development of methods for converting them selectively to other compounds. Contributions to this objective have come from several diverse sources. Biodegradation of petroleum hydrocarbons is a well-known phenomenon and is useful in decontamination of oil leaks and spills. Closely related is utilization of the metabolic apparatus of individual microorganisms to carry out a specific reaction, such as hydroxylation of the C-H bond at the 11 position of a steroid. Many such biochemical oxidations occur by means of a cytochrome P-450 enzyme system in which an oxygen atom is transferred from dioxygen (O2) to an alkane at the iron atom in a heme protein. Synthetic metalloporphyrins and other iron complexes have been studied as model systems for these oxidations. See BIODEGRADATION; CYTOCHROME; PORPHYRIN. James A. Moore Activation by organometallic compounds. Another contribution to the activation of alkanes is from the field of organometallic chemistry. One way that C-H bonds have been cleaved involves the formation of a vacant site on a metal in a low oxidation state that contains electron-donating ligands. The vacant site provides room for the incoming hydrocarbon, and the electron-rich metal readily gives up its electrons in a process called oxidative addition. The net reaction results in the use of two electrons from the C-H bond plus two electrons from the metal to form two metal-carbon bonds, each with two electrons. The vacant site is most easily generated by irradiation of the precursor complex with ultraviolet light, a technique that makes the ligands fall off easily at ambient temperatures. If the irradiation is carried out in a pure hydrocarbon solvent, the metal oxidatively adds to the solvent, resulting in the formation of a complex with metal-hydrogen and metal-carbon bonds. The C-H bond has been completely cleaved, and further chemistry can now be envisioned with the fragments attached to the metal. Reaction (2) shows several specific metal complexes that will undergo this type of oxidative addition chemistry with saturated and unsaturated hydrocarbons such as methane, ethane, hexane, and benzene. In reaction (2), M = rhodium (Rh) or iridium (Ir); Me = methyl group (CH3); Me3P = trimethylphosphine; L = PMe3 or carbonyl (C O); and R = CH3, C2H5 (ethyl), C6H13 (hexyl), or C6H5 (phenyl). With linear alkanes such as hexane, there is a preference of cleavage of the less hindered primary (terminal) C-H bonds at the end of the molecule over the more hindered secondary (internal) C-H bonds of the molecule. In one application based on this reactivity, a complex has been used to convert alkanes into alkenes plus hydrogen by using light energy. The complex chlorocarbonyl bis-(trimethylphosphine) rhodium
459
photolysis
Me3P
M
H
_H
2
R
H
_ CO
H
L
M
L
M
(2) R
H photolysis
OC
M
CO photolysis
−co
RhCl(CO)(PMe3)2
RhCl(PMe3)2
H2
RCH2CH3
CO
(3) RCH2CH2
H H
RhCl(PMe3)2
H R H
RhCl(PMe3)2 H
R
[RhCl(CO)(PMe3)] reacts with alkanes to give many thousands of turnovers of alkenes upon photolysis, as in reactions (3). Beginning with the rhodium complex RhCl(CO)(PMe3)2, the fundamental step involving activation of the C-H bond can be seen in the catalytic cycle following the loss of carbon monoxide (CO). The loss of a hydrogen atom on the carbon that is adjacent to the one bound to the metal is a common reaction in organometallic chemistry, and it leads to the production of the alkene product still bound to the rhodium. The olefin is then released from the complex. Finally, hydrogen is displaced by reaction with CO to regenerate the initial metal complex. Light energy is required in this cycle to remove CO from the metal. See ORGANOMETALLIC COMPOUND. William D. Jones Bibliography. J. A. Davies et al. (eds.), Selective Hydrocarbon Activation, 1990; R. J. Fessenden and J. S. Fessenden, Organic Chemistry, 6th ed., 1998.
Alkene A hydrocarbon that contains a carbon-carbon double bond (C C); dienes or polyenes contain two or more double bonds. The synonym olefin is often used, especially in connection with industrial-scale production and applications. The term unsaturated distinguishes compounds such as fats that contain double bonds from saturated molecules that do not. Structure and names. The carbon atoms joined by a double bond have three planar sp2 orbitals and a p orbital with lobes above and below. The double
RhCl(PMe3)2
460
Alkene volume by catalytic cracking of higher hydrocarbons or dehydrogenation of alkanes during petroleum refining. Alkenes in the C10-20 range are important starting materials in several large-scale industrial applications. These are prepared by controlled oligomerization of ethylene with a triethylaluminum catalyst [reaction (2), where n = 5–10].
W
X C
C Z
Y
(CH3CH2)3AI
Fig. 1. Molecular orbital representation of a carbon double bond; W, X, Y, and Z represent attached groups.
bond is a combination of a single σ bond between the nuclei and a π bond resulting from an electron pair in overlapping p orbitals (Fig. 1). As a result, the attached groups w, x, y, and z are coplanar, and rotation around the C-C axis is possible only by disrupting the π bond. A double bond is shorter (0.133 nanometer) and stronger (610 kilojoules) than a typical C-C bond (0.154 nm, 347 kJ). See CHEMICAL BONDING; MOLECULAR ORBITAL THEORY. The names of alkenes are formed by changing the -ane ending of the corresponding alkane to −ene, with a number to indicate the double-bond location. Since the double bond is a barrier to rotation, an alkene such as 2-butene has two stereoisomeric forms. These geometric isomers are designated as E or trans (similar groups on opposite side) or Z (cis, similar groups on same side), as in the structures below. For the simplest alkenes (C2–C4) the ending H3C
H3C
H C
CH3
C
H
C CH3
H3C C
C H
H
E (trans) 74%
C
(1) H
Z (cis) 23% CH3CH2HC
C
COH
O
–RCO2H
H
C
⌬
O
H
C
CH2
1- Butene 3%
See ISOMERIZATION. Source and preparation. Ethylene (H2C CH2) and the C3–C5 alkenes are produced in very large
C
(3)
R
Ester
hydroxycarbonyl system is the second step in many carbonyl condensation reactions [reaction (4)]. O
OH OH–
O
C H
(4)
CH2C
–H2O
CH3
(2)
C C
CH + H3CC
H
CH2
See CRACKING; DEHYDROGENATION. Several general methods for preparing alkenes involve elimination reactions. Two long-known examples are removal of a hydrogen-halogen group (HX) from a halide with strong base, and loss of water from alcohols. Both reactions can often give more than one alkene, and acid-catalyzed dehydration may lead to rearrangement. In the latter case, conversion of the alcohol to an ester and pyrolysis may be preferable [reaction (3)]. Dehydration of a β-
Z-2-Butene
CH3 C
H3C
CH3CH2(CH2CH2)n –2HC
O
-ylene is often used. Two other common names are vinyl for the group HC CH2 and allyl for the group CH2HC CH2. See ALKANE. An alkene can undergo isomerization by shifts of the double bond, such as by treatment with an acid catalyst. Isomers with the double bond in different positions or geometries have different internal energies and therefore different stabilities. In most cases a trans isomer is favored over the cis, because larger groups are farther apart. Moreover, isomers in which the double bond is in an internal position on the chain are usually more stable than those with the double bond at the end. For the three straightchain butene isomers, the equilibrium composition is shown in notation (1). H
CH2
C
H
E-2-Butene
n HC2
HC
O CHC
A widely used preparative method for alkenes is the Wittig reaction. In this process a phosphorane adds to a carbonyl compound, and the double bond is formed by loss of triphenylphosphine oxide from a phosphetane [reaction (5)]. Another approach to RCH2Br + PPh3
RCH2PPh3 Phosphonium salt
RHC
PPh3
Phosphorane C
O
(5) Ph3P O + C CHR
O
PPh3
O
PPh3
C
CHR
C
CHR
Phosphetane
the formation of a double bond is reduction of a triple bond. See ALKYNE; ORGANOPHOSPHORUS COMPOUND. Chemical activity. The double bond is a highly reactive and versatile functional group. Because of the readily accessible electrons in the π bond, alkenes react with a variety of Lewis acids, or electrophiles.
Alkene A characteristic reaction is addition, with formation of two new σ bonds in place of the π bond. Hydrogenation. Addition of hydrogen occurs on the surface of a finely divided metal catalyst, usually nickel, palladium, or platinum. This reaction is useful for several purposes; one application is the partial hydrogenation of unsaturated vegetable oils to remove some of the double bonds and raise the melting temperature. See FAT AND OIL (FOOD). Proton acids. Several types of addition products can arise from the reaction of alkenes with acids. In a simple case hydrochloric acid (HCl) or hydrobromic acid (HBr) add to propene to give the 2-halopropane; with aqueous sulfuric acid the elements of water (H2O) add, and 2-propanol is obtained [reaction (6),
of the original double bond. When the reaction is carried out in an aqueous or alcoholic solution of the halogen, water or alcohol is the nucleophile in the second step and the product is the bromo alcohol or ether. H
H
Br Br2
H
C
C R
R
CH3HC
[CH3CHCH3]
CH3CHCH3
(6) +
[CH3CH2CH2]
where X= chlorine (Cl), bromine (Br), or hydroxyl (OH)]. These additions occur by way of an intermediate carbenium ion that rapidly reacts with a nucleophile such as halide ion or water. The structure of the product depends on the fact that the cation with the larger number of alkyl groups attached to the electron-deficient carbon is more stable and more easily formed. See REACTIVE INTERMEDIATES. The carbenium ion formed from an alkene may undergo other reactions as well. Under suitable conditions it can react as an electrophile with another molecule of the alkene to give branched dimers, trimers, and so forth [reaction (7), where R= CH3 RCH
CH2
H+
RCHCH3
RHC
CH2
R
C
R
OH
H
H
C
C
R
OR′
H C
R
CH2
O
O2, Ag, or
H
C R
RCO3H
C
C
R
H
(9)
R Epoxide
See EPOXIDE. Alkenes can also be oxidized to diols. This reaction takes place when a transition-metal oxide such as permanganate (MnO4) or osmium tetroxide adds to the double bond. A cyclic intermediate results, and this is hydrolyzed to give the diol [reaction (10)]. (OsO2)
+
CHR
–
H+
dimer
(7)
Dimeric cation RHC
R'
(8) Br
H
C
C
+ OsO4
O
O
C
C
H2O
Cyclic intermediate
CH2
[trimeric cation]+
C
– H+
trimer
– H+
or C2H5 and R = C9H18 or C12H24]. An alkene with acid can also be used to alkylate a benzene ring in a variant of the Friedel-Crafts reaction. Both processes are used in the production of alkylbenzenes for detergents and other uses. See FRIEDEL-CRAFTS REACTION. Halogenation. Bromine or chlorine adds very rapidly to alkenes; bleaching of the red color of Br2 has long been recognized as a test for a double bond. The reaction of ethylene with chlorine to give 1,2dichloroethane is the first step in the manufacture of vinyl chloride. In the addition of halogen to a double bond [shown in reaction (8) for Br2], a cyclic cation is initially formed. In the second step, bromide ion attacks this intermediate at a C-Br bond. The result is anti addition, with the Br atoms bonded to opposite sides
R Br
See HALOGENATED HYDROCARBON. Oxidation. Addition of oxygen to a double bond gives an epoxide. This reaction is carried out on an industrial scale by passing an alkene and air over a silver (Ag) catalyst. In smaller-scale work, epoxides can be obtained by treating an alkene or other unsaturated compound with a peroxide or peracid [reaction (9)]. H
H C
R′OH
Br H
C R
H2O
CH3 RCH
H
H
C
X–
+
CH2
Br
_
Br
R R Cyclic cation
X H+
C
⊕
OH
OH
C
C
(10)
Diol Hydroboration. Alkenes combine readily with borane (BH3). The resulting alkyl boranes are highly reactive compounds and are useful and versatile reagents in organic synthesis. See HYDROBORATION. Free radicals. Double bonds undergo addition reactions by attack of free radicals, X·. Some of the important radicals are obtained from HX compounds such as hydrogen bromide (HBr), hydrogen sulfide (H2S), a thiol (HSR), and chloroform (HCCl3). The process is initiated by homolysis of HX, often with air [reaction (11), where X=HBr, HSR, or HCCl3]. A chain reaction then begins (propagation), with X· adding to the π bond [reaction (12)]. Subsequent reaction of
R
461
462
Alkene the intermediate radical with another HX molecule gives a product and a new X· radical to continue the chain [reaction (13)]. The direction of addition Initiation:
(11)
X + HO2
HX
CH2
H
H2O
CHOH
Pd+
CH2
CH2
RCHCH2X
CH2 + X
O
Cu2+
H
CHOH
(13)
is due to stabilization of the electron-deficient radical by alkyl substituents, just as in the case of a carbenium ion [reaction (6)]. When X = Br, it will be noted that addition by the radical mechanism [reactions (12) and (13)] occurs in the direction opposite to that in the ionic mechanism. The observation and clarification of this reversal of the addition of HBr was a milestone in the understanding of how reactions occur. See CHAIN REACTION (CHEMISTRY). The free-radical intermediate in reaction (12) can react with HX as in reaction (13), but it can also react with another alkene molecule, or several molecules, before being capped. This possibility is shown in reaction scheme (14) for the free-radical addition of CI + CCI3
H2C
CH2
CCI3CH2CH2 n –1 H2C
CCI3(CH2CH2)nCI
CCI4
CH2
carbon tetrachloride (CCl4) to ethylene. This reaction results in a mixture of products; the main products have the following yields; n = 1, 9%; n = 2, 57%; and n = 3, 24%. The overall process is called telomerization. See FREE RADICAL. Reactions of metal complexes. Among the most important properties of alkenes is the formation of π complexes by donation of an electron pair to the coordination shell of a transition metal. This is a very general reaction, and the complexes are central to some of the most important chemistry, including industrial processes, involving alkenes. Chief among these is the production of polyethylene and polypropylene by polymerization of the alkenes with Ziegler-Natta catalysts, in which monomer units are incorporated into the chain at a coordination site on a transition metal such as titanium(III). See COORDINATION COMPLEXES; ORGANOMETALLIC COMPOUND; POLYMERIZATION; STEREOSPECIFIC CATALYST. In complex formation, electron density shifts from the π bond to the metal, and the coordinated alkene is therefore susceptible to addition of nucleophiles such as an OH or NH group. Palladium (Pd) is a particularly versatile metal in several reactions. In the Wacker process for the production of acetaldehyde, ethylene is complexed with PdII [reaction (15)]. After addition of water, the complex undergoes βelimination of enol, and palladium is thereby reduced
Acetaldehyde
to Pd0. It is reoxidized to PdII by copper ion (Cu2+), which is in turn reoxidized by oxygen (air). The net reaction is thus a two-electron oxidation of ethylene, with both palladium and copper required in only catalytic amounts. Analogous reactions in which alcohol or acetic acid are used rather than water produce vinyl ether and vinyl acetate, respectively. Extension of this reaction to higher 1-alkenes provides a convenient and practical method to prepare methyl ketones [reaction (16)]. O Pd ⌸, air
RHC
CH2
H2O
RCCH3
(16)
Methyl ketone
Another example of the utility of palladium π complexes is the substitution of an aryl (Ar) or vinyl group for a vinyl hydrogen in an alkene [reaction (17)]. In this reaction the aryl group is transferred
(14)
CCI3(CH2CH2)n
CH3CH
CH2 Enol
(12)
RCH2CH2X + X
(15)
β-elimination
/2 O2,
Pd
RCH CH2X + HX
hν
Pd
Pd0 + H+
Propagation: RHC
CH2 CH2
1
O2
CCI4
Pd2+ +
Ar H ArX + RC
Pd ⌸
CR2
RC
CR2 + HX
(17)
from the metal to the complexed alkene, followed by β-elimination of the product. Hydroformylation is a major industrial reaction in which an alkene is converted to an aldehyde containing an additional carbon. The alkene is complexed to cobalt (Co) in the form of HCo(CO3)3. Transfer of hydrogen from the metal is followed by insertion of carbon monoxide (CO) and finally reduction of a CoCO bond [reaction (18)]. See HYDROFORMYLATION. A different type of organometallic chemistry underlies a process known as olefin metathesis, in which the groups attached to the double bonds in two alkenes are scrambled. The reaction proceeds by way of metal carbene complexes of molybdenum or tungsten (W), and the reversible formation of metallacyclobutanes [reaction (19); L = CO]. In this way an alkene with a double bond in the center of a chain can be converted by reaction with ethylene to two new terminal alkenes. Allylic compounds. The position adjacent to a double bond is termed allylic, and a group or atom at that position is strongly affected by the π bond. Thus in a radical or an ionic intermediate, the electron deficiency or excess is delocalized by conjugation with the p orbitals (Fig. 2). As a consequence, these species are lower in energy and more easily formed than those in a saturated molecule. For this reason, several reactions of alkenes occur by
Alkene
463
R H CH CH2
HCo(CO)3 + RHC
(CO)3Co
(CO)3Co CH2
CH2CH2R
(18)
CO
O H2
RCH2CH2CH
Fig. 2. Molecular orbital representation of a reactive intermediate, with delocalization by conjugation with the p orbitals.
R1HC
substitution of the allylic position rather than addition at the double bond. Examples are the catalytic oxidation of propene to acrylic acid or acrylonitrile, or free-radical chlorination to allyl chloride [reaction (20)]. See DELOCALIZATION. Reactions such as hydrolysis of an allylic type chloride often occur with rearrangement, since the allylic carbenium ion can be solvated at either end of the allylic system. An example is 1-chloro-2-butene, which gives mainly 1-buten-3-ol and a small amount of 2-buten-1-ol. The other chloride gives the same mixture of alcohols [reaction (21), where the positive charge and the π bond in the intermediate are delocalized over two carbon atoms]. Allylic compounds readily form π complexes with palladium, nickel, and other metals. These π-allyl complexes, in which the metal is complexed to three carbons, provide selective activation of the allylic system for C-C bond formation with a nucleophilic carbon, as in reaction (22) where Nu represents the nucleophile and X = Cl, RO, or ROC = O. Dienes. Two double bonds, for example in a pentadiene, can exist in three arrangements: cumulative or allenic, conjugated, and nonconjugated, as in 2,3-, 1,3-, and 1,4-pentadiene, respectively. The 2,3-isomer is an allene, for example, 2,3-pentadiene (CH3HC C CHCH3); 1,3pentadiene (H2C CH-CH CHCH3) is conjugated; 1,4-pentadiene (H2C CHCH2CH CH2) is nonconjugated. In 2,3-pentadiene the central carbon has sp hybrid orbitals and is susceptible to attack by either electrophilic or nucleophilic reagents. A disubstituted allene such as 2,3-pentadiene is a chiral molecule and exists in enantiomeric forms (Fig. 3). The characteristic feature of a conjugated diene is 1,4-addition. In the reaction of 1,3-pentadiene with bromine, electrophilic attack gives an allylic cation which can lead to either 1,4- or 1,2-addition [reaction (23)]. The major product generally results from the former. Another important aspect of conjugated diene chemistry is a group of reactions that take place simply by shifts of σ and π bonds. These include ring opening of cyclohexadienes, Diels-Alder cyH H3C
H CH3
H
mirror
H3C
H CH3
Fig. 3. Enantiomeric forms of the disubstituted allene 2,3-pentadiene. The bold dot represents the central carbon atom.
+
L4W
(CO)3CoCOCH2CH2R
CHR2
R1CH
CHR2
R1CH
CH2
L4W
CH2
L4W
CHR2
+
CH2
(19) R1HC
CH2
H C
[O], NH3
CH3
CH2
L4W
CH2
H2C
[O]
H2C
R1CH
H2C
CHCO2H Acrylic acid
H2C
CHCN Acrylonitrile
H2C
CHCH2CI Allyl chloride
CI2, h
CH2
(20)
CI CH3HC
CHCH2CI
␦+ _ _ _ _ _ _ _ ␦+ CH3CH CH CH2
H2O
H2O
1-Chloro-2-butene
CH3CHHC
CH2
3-Chloro-1-butene H2O
(21)
OH
C
X
C
Major:
CH3CHC H +
Minor:
CH3HC
C
Pd ⌸
C
C
1--Buten-3-ol
CH2
2--Buten-1-ol
CHCH2OH
C
Nuc:
C
C
Nuc C
(22)
Pd
H2C
H C
H C
+
__________
Br2
CHCH3
BrH2C
CH CH CHCH3 Allylic cation
1,4-addition
BrCH2HC
CHCHBrCH3
1,2-addition
BrCH2CHBrHC
CHCH3
cloadditions, and double-bond (sigmatropic) rearrangements, referred to collectively as pericyclic reactions. See CONJUGATION AND HYPERCONJUGATION; DIELS-ALDER REACTION; PERICYCLIC REACTION; WOODWARD-HOFFMANN RULE. In nonconjugated dienes the double bonds can react independently, but with a 1,5- or 1,6-diene electrophilic attack at one double bond may be
(23)
464
Alkylation (petroleum) accompanied by cyclization [reaction (24), where E + H
(24)
H
E+
E
E
represents the electrophile]. This process frequently occurs in the biosynthesis of steroids and terpenoids. Dienes form complexes of several kinds with transition metals. With PdII salts, butadiene forms π-allyl complexes. Complexes of iron with conjugated dienes are coordinated to four atoms. Typical compounds are linear or cyclic diene–iron tricarbonyl complexes. The elusive cyclobutadiene forms stable complexes with iron and nickel; the bonding is delocalized to achieve an 18-electron count around the metal, as in metallocenes. See METALLOCENES. Metal complexes are involved in the oligomerization of dienes to cyclic dimers and trimers. Polymerization of 1,3-dienes occurs by 1,4-addition and can lead to either cis- or trans-linked chains [reaction (25),where R CH3, the 1,3-diene is isoprene, H
R C
C
CH2 CH2
CH
CR
CH2
cis
CH2
(25) H
CH2 C
C
CH2
R
trans
and where R Cl, the 1,3-diene is chloroprene]. The stereochemistry can be controlled by the choice of catalysts. These polymers represent the structures of elastomeric materials. Cis-1,4-(poly)isoprene obtained by polymerization with R3Al-TiCl4 catalysts is the basic unit of natural rubber from Hevea species. The trans polyisoprene is obtained with other R3Al catalysts and is the repeating unit of gutta-percha. Addition of small amounts of cis-1,4(poly)butadiene significantly improves the wear resistance of natural rubber. Neoprene, the polymer of 2-chlorobutadiene (chloroprene), has several advantageous features. See RUBBER. Polyenes. Compounds with three or more double bonds are found in numerous naturally occurring sources. The long-chain (C12–C20) carboxylic acids present as esters in fats and other lipids typically contain two or three nonconjugated double bonds. The most conspicuous polyenes are carotenoids, the red-yellow pigments that occur abundantly in plants and many other organisms. These compounds are typically C40 molecules with an extended system of conjugated trans double bonds. Vitamin A is a pentaene derived in the organism by cleavage of βcarotene. Reversible cis-trans isomerization of one double bond plays a central role in color vision. See CAROTENOID; VITAMIN A.
A very long polyene chain is obtained by polymerization of acetylene with an organometallic catalyst. This material can be treated to give an electrically conducting polymer. See ORGANIC CONDUCTOR. James A. Moore Bibliography. A. Streitwieser, C. Heathcock, and E. M. Kosower, Introduction to Organic Chemistry, 4th ed. 1992.
Alkylation (petroleum) In the petroleum industry, a chemical process in which an olefin (ethylene, propylene, and so forth) and a hydrocarbon, usually 2-methylpropane, are combined to produce a higher-molecular-weight and higher-carbon-number product. The product has a higher octane rating and is used to improve the quality of gasoline-range fuels. The process was originally developed during World War II to produce highoctane aviation gasoline. Its current main application is in the production of unleaded automotive gasoline. See GASOLINE; OCTANE NUMBER. Reactions. The alkylation reaction is initiated by the addition of a proton (H+) to the olefin [reaction (1)]. The protonated olefin (carbonium reactions) then reacts with the isobutane by abstraction of a proton from the isobutane to produce the t-butyl carbonium reactions [reaction (2)]. Reaction of this tertiary carbonium reactions with the olefin proceeds by combination of the two species [reaction (3)] to produce a more complex six-carbon carbonium reation which yields a stabilized product by abstraction of a proton from another molecule of isobutane [reaction (4)]. The reaction progresses to more product by reaction of the t-butyl and carbonium reactions, as already shown in reaction (3). See REACTIVE INTERMEDIATES. H2 C
CH2 +
CH2 + H+ → H3 C
(1)
CH3 H3C
CH + H3C
+
CH2
CH3
CH3 H3C
C+ + H3C
CH3
(2)
CH3 CH3 H3C
C+ + H2C
CH2
CH3
CH3 H H3C
C
C
+
CH2
(3)
CH3 H
Processes. In actual refinery practice, the feedstock for the alkylation process is isobutane (or an isobutane-enriched stream) recovered from refinery
Alkyne CH3
CH3 H H3C
C
C
+
CH2 + H3C
R
CH
CH3 H
H3C
C
C
CH3 H
R
C
Fig. 1. Geometry of a linear arrangement of triply bonded carbon atoms and two attached groups.
CH3 CH3 H
C
CH3 CH3 + H3C
C+
(4)
CH3
Alkylate
gases or produced by isobutane isomerization. The olefin for the feedstock is derived from the gas production of a catalytic cracker. See CRACKING; ISOMERIZATION. The acid catalyst is sulfuric acid (H2SO4), hydrofluoric acid (HF), or aluminum chloride (AlCl3). When sulfuric acid is employed, the reaction must be maintained at 35–45◦F (2–8◦C) to reduce unnecessary side reactions; with hydrogen fluoride and with aluminum chloride, refrigeration is not necessary, and temperatures up to 115◦F (45◦C) may be employed. Sulfuric acid is used with propylene and highercarbon-number olefins but not with ethylene because of the tendency to produce the acid-wasting ethyl hydrogen sulfate (C2H5HSO4). Hydrogen fluoride is more generally used, and has the advantage of being more readily separated from the products because of a boiling point below that of sulfuric acid. Aluminum chloride is used to a lesser extent than both of the acids and requires injection of water to produce the hydrogen chloride promoter. See CATALYSIS. Thermal alkylation is a noncatalyzed process that requires temperatures on the order of 950±25◦F (about 510◦C) and pressures in the range 3000– 8000 lb/in.2 (18–55 megapascals). The reaction is believed to proceed by way of a (neutral) free-radical intermediate and is much less specific than the acidcatalyzed process. The starting hydrocarbon is usually isobutane, and the olefin is ethylene, propylene, or higher olefins such as butenes. The nature of the reactants influences the quality of the alkylate. For example, alkylation of 1-butene with isobutane gives an alkylate with a research octane rating of about 93, whereas using 2-butene as the olefin yields an alkylate having a research octane number of about 98. See ALKANE; ALKENE; FREE RADICAL; PETROLEUM PROCESSING AND REFINING. James G. Speight Bibliography. S. Glasstone, Energy Deskbook, 1983; J. McKetta, Petroleum Processing Handbook, 1992; J. G. Speight, The Chemistry and Technology of Petroleum, 3d ed., 1999.
Alkyne A hydrocarbon that contains a triple carbon bond, C C . The first compound in the series is acetylene. Names of alkynes are formed by changing the -ane ending of the corresponding alkane to -yne and numbering as needed, for example, 2-pentyne (struc-
ture 1), the 2 indicating that the triple bond is at the second carbon atom of the chain. Some compounds can be named as a substituted acetylene, for example, diphenylacetylene (2); or the ethynyl group ( C CH) can be named as a substituent, for example, 1-ethynylcyclopentanol (3).
CH3CH2C
C6H5C
CCH3
CC6H5
(2)
(1) OH C
CH
(3) See ACETYLENE; ALKANE. Structure and properties. A triple bond consists of a σ bond and two π bonds between two sp carbon atoms. The carbons are held closer together (0.121 nanometer) and more strongly (836 kilojoules/mole) than those in a double bond. The π bonds are orthogonal, and the electron density is distributed in a cylindrical sheath around the C-C axis. The resulting geometry is a linear arrangement of the triply bonded carbons and the two attached atoms (Fig. 1). See MOLECULAR ORBITAL THEORY. The energy difference between a 1-alkyne and 2-alkyne (20 kJ) is larger than that between 1- and 2-alkenes, and a 1-alkyne can be completely isomerized with a base catalyst [reaction (1)]. An intermeRCH2 C
CH − − RCH
C
CH2 −→ RC
CH3
diate in this reaction is the allene isomer, which is only slightly lower in energy than the 1-alkyne. Cycloalkynes can be prepared, but the requirement of a linear four-atom system for a triple bond places a limit on the minimum size of a cycloalkane ring. Cyclooctyne, with an eight-membered ring, can be isolated; seven- and six-membered cyclic alkynes can be trapped as transient intermediates. In cycloalkynes with medium rings (8–11 carbons) the amount of allene isomer is substantial, since only three atoms are collinear and the cyclic allene in this case is significantly more stable than the cycloalkyne. See ISOMERIZATION; REACTIVE INTERMEDIATES. A distinctive property of acetylene and other 1-alkynes is the acidity of the C C H bond. Because the electrons in an sp carbon orbital are closer to carbon than those in sp2 or sp3 orbitals, the anion C C:− is stabilized. The acidities of a few relevant compounds are given in the table. The terminal C-H
465
466
Alkyne [reaction (5)] or cyclobutanedione (7) [reaction
Comparison of acidity values of several compounds
Compound
Acidity, pK A (⫺log KA )
Ethane (H3 C CH3 ) Ethylene (H2 C CH2 ) Ammonia (NH3 ) Acetylene (HC CH) Ethanol (CH3 CH2 OH)
50 44 35 25 16
RC
CH + NaNH2
Alkyne
RC
CNa + NH3
–CO
R
O
in acetylene (and other alkynes) is far more acidic (about 1020–1025 times) than hydrogen in alkenes or alkanes. Hydrogen can be removed from an alkyne by reaction with sodium amide, a reaction that is conveniently carried out in liquid ammonia [NH3; reaction (2)]. Many other bases, such as alkyllithiums, can NH3
O
(2)
Sodium acetylide
be used to prepare the acetylide anion ( C C: ). See PK. Preparation. One general method for introducing a triple bond into a carbon chain is by elimination reactions, for example, removal of HX ( X = halogen) from a vinyl halide [H2C CHX] with very strong base, or removal of two HBr molecules from a dibromide [reaction (3)].
Cl
RC
CBr2
HBrC
CHBr
RC
CR
R
(C
(5)
CH
(7)
Cl O
O
RC
C
C
CR
⌬
C)3
R
(6)
(6)]. The latter reaction can be used to build up chains containing several triple bonds. A second general approach for introducing a triple bond is using reactions involving an acetylide anion in either a substitution or addition reaction. The acetylide anion is an excellent nucleophile, similar in reactivity to cyanide ion (N C: ). Any alkyl halide or sulfonate ester that is a practical substrate for substitution can be converted to an alkyne. With a dihalide (8) and excess acetylide, a diacetylene (9) is available [reaction (7)]. Br(CH2)nBr + 2HC
CNa
Sodium acetylide
(8) H2C
R
(6) O
⌬
HC
– HBr
C(CH2)nC
CH
(7)
(9) – HBr
HC
CBr
C
C
(3)
Another starting point for elimination is an enol ester such as a trifluoromethylsulfonate (4) or an enol phosphate (5), as in reactions (4). The latter O (CF3SO2)2O
CH2C Carbonyl compound OSO2CF3 CH
C
(4) OPO3H CH
C
C
(4)
Addition of acetylides to aldehydes, ketones, or epoxides leads to the corresponding alkynols. An acetylide can be acylated with an acid chloride or carbonated to give the ketone or carboxylic acid, respectively. Addition reactions. In general, electrophilic additions are slower with triple bonds than double bonds because π electrons in the latter are more accessible. However, several addition reactions of alkynes are useful. Hydrogenation of a triple bond can be carried out with a palladium (Pd) catalyst that is deactivated [“poisoned” with the addition of lead (Pb)] to permit selective reduction to the cis-alkene. Reduction to the trans-alkene can be accomplished with sodium metal (Na) in ammonia (NH3) containing alcohol [reaction (8)]. R
C
(5) has been suggested as the precursor of the triple bonds in naturally occurring acetylenes. Alkynes are also obtained by thermal elimination of neutral molecules, as in the preparation of cycloalkynes by oxidizing the hydrazones of 1,2diketones. Another reaction is elimination of carbon monoxide (CO) from a cyclopropenone (6)
R' C
H
2Na
RC H
cis
H2
C
C
R'
Pd(Pb)
NH3,ROH
R
H C
C
H
R'
(8)
trans
Addition of water to a terminal triple bond occurs in the presence of a mercuric (HgII) salt, leading
Alkyne initially to an enol and thence to a ketone. Reaction of the alkyne with borane (B2H6) gives a vinyl borane that can be oxidized in the presence of water (H2O) and hydrogen peroxide (H2O2) to the isomeric enol of the aldehyde [reactions (9)]. Dialkyl aluminum
are usually good, and a variety of functional groups, including additional triple bonds, can be present. Oxidative coupling of terminal dialkynes leads to macrocyclic polyacetylenes [reaction (11)]. This is C
OH H2O
RC
RC CH2 Enol
CH
HgII
RHC CHB Vinyl borane
(9)
H2O2
O RCH2CH Aldehyde
RCCH3 Ketone
O2
(CH2)n C
H2O
O
CH
(CH2)n
B2H6
RHC CHOH Enol
reagents add to alkynes. The alkenylaluminum products can be converted stereoselectively to alkenes or to alkenyl halides or allylic alcohols. Cyclization reactions. In early work benzene was observed among the products formed on heating acetylene. Substituted acetylenes are quite stable thermally, but metal-catalyzed thermal trimerization gives trisubstituted benzenes, with the 1,2,4-isomer greatly predominating over the 1,3,5-isomer [reaction (10)].
CH
(CH2)n
Cu2+
(11)
a very direct way to prepare large-ring compounds, although rings with different numbers of diyne units and also polymeric products are formed. See MACROCYCLIC COMPOUND. Polyynes occur in many species of higher plants, particularly in the families Compositae, Umbelliferae, and Asteraceae as well as in Basidiomycetes and other fungi. Many hundreds of compounds have been isolated and characterized; nearly all of these are structurally and biogenetically related to fatty acids with unbranched C10–C20 chains. A number of compounds are derived from a series of C13 hydrocarbons that contain a combination of six triple and double bonds. A common feature is the occurrence of thiophenes arising from addition of the biochemical equivalent of hydrogen sulfide (H2S) to two adjacent triple bonds [reaction (12)]. Many other functional H2S
metal catalyst
3(CH3)3CC
467
CH
CH3HC
CH(C
C)4HC
CH2 enzyme
C(CH3)3
C(CH3)3
CH3HC
CH
S
C(CH3)3
(10) (CH3)3C
C(CH3)3
C(CH3)3
This bias for the 1,2,4 pattern indicates a mechanism that is not a simple head-to-tail cyclization, and may involve some type of metal-oriented dimeric complex. The reaction is a route to benzenes with adjacent bulky groups. The formation of cyclooctatetraene in high yield by reaction of acetylene in the presence of a nickel (NiII) catalyst was a seminal finding in organometallic chemistry. A variety of other cyclizations of alkynemetal complexes have been observed. In some of these reactions carbon monoxide molecules are also involved, leading to cyclopentadienones and tropones. The cocylization of an alkyne with an alkene and carbon monoxide, catalyzed by dicobalt octacarbonyl [Co2(CO)8], is a general and flexible method for cyclopentenone preparation. See ORGANOMETALLIC COMPOUND. Diynes and polyynes. 1-Alkynes can be coupled by mild oxidation to the conjugated diacetylene. The reaction is carried out by air oxidation of the cuprous (CuI) acetylide or oxidation of the alkyne with cupric (CuII) ion in the presence of an amine. For coupling of two different alkynes, the reaction is carried out with the copper derivative of one alkyne and a different 1-bromoalkyne. The yields in these couplings
HC 2
CH2 (12)
groups such as hydroxyl, epoxy, furan, and pyran rings are also present. A group of antibiotics that are potent inhibitors of tumor growth have as a common feature a macrocyclic enediyne system. The mechanism of action of these compounds involves rearrangement of the multiple bonds to a benzene diradical which then interferes with deoxyribonucleic acid (DNA) synthesis. See ANTIBIOTIC. Compounds containing an extended sequence of conjugated triple bonds with the general formula (C C)n have been detected in interstellar matter as well as in plant sources. Polyyne chains with up to 16 triple bonds (Fig. 2a), and also molecules such
(C
(a)
C)
16
(b)
(c)
Fig. 2. Compounds containing an extended sequence of conjugated triple carbon bonds: (a) polyyne chain with 16 triple bonds; (b) hexa(butadiynyl)benzene; and (c) tetrabutadiynyl ethylene.
468
Allantois as hexa(butadiynyl)benzene (Fig. 2b) and tetrabutadiynyl ethylene (Fig. 2c), have been prepared. These compounds are stable with a bulky group such as trialkylsilyl at the ends of the polyyne chains; but with terminal hydrogen, they are extremely reactive and may ignite in air or explode. Compounds with extended polyyne systems are of interest as precursors of carbon rods or two-dimensional carbon networks, both of which have potential as new materials. James A. Moore Bibliography. L. Brandsma and H. D. Verkruijsse, Syntheses of Acetylenes, Allenes and Cumulenes, 1981; F. Diedrich and Y. Rubin, Synthetic approaches to molecular carbon allotropes, Angew. Chem. Int. Ed., 31:1101–1123, 1992; H. G. Viehe (ed.), Chemistry of Acetylenes, 1969.
Allantois A fluid-filled sac- or sausagelike, extraembryonic membrane lying between the outer chorion and the inner amnion and yolk sac of the embryos of reptiles, birds, and mammals. The allantois eventually fills up the space of the extraembryonic coelom in most of these animals. It is composed of an inner layer of endoderm cells, continuous with the endoderm of the embryonic gut, or digestive tract, and an outer layer of mesoderm, continuous with the splanchnic mesoderm of the embryo. It arises as an outpouching of the ventral floor of the hindgut and dilates like a filling balloon into a large allantoic sac which spreads throughout the extraembryonic coelom. The allantois remains connected to the hindgut by a narrower allantoic stalk which runs through the umbilical cord. See AMNION; CHORION; GERM LAYERS. The allantois eventually fuses with the overlying chorion to form the compound chorioallantois, which lies just below the shell membranes in reptiles and birds. The chorioallantois is supplied with an extensive network of blood vessels and serves as an important respiratory and excretory organ for gaseous interchange. The allantoic cavity also serves as a reservoir for kidney wastes in some mammals, in reptiles, and in birds. In the latter two groups the allantois assists in the absorption of albumin. In some mammals, including humans, the allantois is vestigial and may regress, yet the homologous blood vessels persist as the important umbilical arteries and veins connecting the embryo with the placenta. See FETAL MEMBRANE; PLACENTATION. Nelson T. Spratt, Jr. Bibliography. B. I. Balinsky, Introduction to Embryology, 5th ed., 1981.
Allele Any of a number of alternative forms of a gene. Allele is a contraction of allelomorph, a term which W. Bateson used to designate one of the alternative forms of a unit showing mendelian segregation. New alleles arise from existing ones by mutation. The diversity of alleles produced in this way is the basis for hered-
itary variation and evolution. See GENE; MENDELISM; MUTATION.
The genetic material is a coded set of instructions in chemical form (deoxyribonucleic acid) for making a living cell. Each gene is like an essential word in the code. The different alleles of a given gene determine the degree to which the specific hereditary characteristic controlled by that gene is manifested. The particular allele which causes that characteristic to be expressed in a normal fashion is often referred to as the wild-type allele. Mutations of the wild-type allele result in mutant alleles, whose functioning in the development of the organism is generally impaired relative to that of the wild-type allele. In this analogy, the wild-type allele corresponds to the correct word in the genetic code and the mutant alleles correspond to mistakes in the spelling of that word. See DEOXYRIBONUCLEIC ACID (DNA); GENETIC CODE. An allele occupies a fixed position or locus in the chromosome. In the body cells of most higher organisms, including humans, there are two chromosomes of each kind and hence two alleles of each kind of gene, except for the sex chromosomes. Such organisms and their somatic cells are said to carry a diploid complement of alleles. A diploid individual is homozygous if the same allele is present twice, or heterozygous if two different alleles are present. Let A and a represent a pair of alleles of a given gene; then A/A and a/a are the genetic constitutions or genotypes of the two possible homozygotes, while A/a is the genotype of the heterozygote. Usually the appearance or phenotype of the A/a individuals resembles that of the A/A type; A is then said to be the dominant allele and a the recessive allele. In the case of the sex chromosomes, one sex (usually the male in most higher animals, with the exception of birds) has only one X chromosome, and the Y lacks almost all of the genes in X. The male thus carries only one dose of X-linked genes and is said to be hemizygous for alleles carried on his X chromosome. As a result, if a male inherits a recessive mutant allele such as color blindness on his X chromosome, he expresses color blindness because he lacks the wild-type allele on his Y chromosome. See CHROMOSOME; SEX-LINKED INHERITANCE. The wild-type allele of a gene is symbolized by attaching a superscript plus sign to the gene symbol; for example, a+ is the wild-type allele of the recessive mutant gene a. When no confusion arises, the symbol + is used instead of a+. This means that a/+ is equivalent to a/a+. In some cases a mutant is dominant to the wild-type allele, and then the mutant symbol is usually capitalized. In a population of diploid individuals, it is possible to have more than two alleles of a given gene. The aggregate of such alleles is called a multiple allelic series. Since genes are linear sequences of hundreds or even thousands of nucleotide base pairs, the potential number of alleles of a given gene which can arise by base substitution alone is enormous. For example, a gene coding for a protein with only 100 amino acids would be considered a very small gene, and would have 300 base pairs in its coding region; yet, since each base pair can exist in four different chemical
Allelopathy forms, 4300 possible mutant alleles could arise within the coding region alone, or approximately 10180 alleles—a number vastly larger than the estimated number of particles in the known universe. Actually, many other types of mutant alleles of a gene do occur, including duplications, deletions, and inversions of DNA within the noncoding as well as coding regions of the gene. Indeed, many spontaneous mutations in at least some higher organisms, such as yeast, maize, and Drosophila, turn out to be insertions of large segments of DNA into either the noncoding or coding regions. These insertions, known as transposons, appear to be virallike elements that have infected the germ line. See TRANSPOSONS. Mutations of independent origin that involve changes at the same nucleotide positions in a gene are known as homoalleles, while mutations that involve changes at different nucleotide positions are called heteroalleles. Recombination between alleles of the same gene does not occur or tends to be extremely rare in higher organisms. Also, recombination within the gene in such cases tends to be different in outward properties from ordinary crossing over and is known as gene conversion. See RECOMBINATION (GENETICS). It is often difficult to distinguish operationally between a series of multiple alleles of a single gene and a cluster of closely linked genes with similar effects. Such a cluster, originally known as a pseudoallelic series, is now termed a gene complex, gene cluster, or multigene family, except that, in the latter case, the genes of the family are not necessarily closely linked. In bacteria, gene clusters are usually coordinately regulated and are termed operons. Gene cloning has verified the existence of gene clusters in higher organisms as well as in bacteria. See BACTERIAL GENETICS; COMPLEMENTATION (GENETICS); GENETICS; OPERON. Edward B. Lewis
Allelopathy The biochemical interactions among all types of plants, including microorganisms. The term is usually interpreted as the detrimental influence of one plant upon another but is used more and more, as intended originally, to encompass both detrimental and beneficial interactions. Forms and occurrence. At least two forms of allelopathy are distinguished: (1) the production and release of an allelochemical by one species inhibiting the growth of only other adjacent species, which may confer competitive advantage for the allelopathic species; and (2) autoallelopathy, in which both the species producing the allelochemical and unrelated species are indiscriminately affected. The term allelopathy, frequently restricted to interactions among higher plants, is now applied to interactions among plants from all divisions, including algae. Even interactions between plants and herbivorous insects or nematodes in which plant substances attract, repel, deter, or retard the growth of attacking insects or nematodes are considered to be allelopathic.
Interactions between soil microorganisms and plants are important in allelopathy. Fungi and bacteria may produce and release inhibitors or promoters. Through the secretion of siderophores they can immobilize iron, which may deprive harmful rhizobacteria and pathogens of this important trace element. Siderophores from pseudomonad bacteria, however, may reduce iron uptake by plants, whereas siderophores from agrobacteria may enhance it. Some bacteria enhance plant growth through fixing nitrogen, others through providing phosphorus. The activity of nitrogen-fixing bacteria may be affected by allelochemicals, and this effect in turn may influence ecological patterns. Allelochemicals, released from plant litter by microbial action, are often also removed from the rhizosphere by microorganisms. The rhizosphere must be considered the main site for allelopathic interactions. See NITROGEN FIXATION; RHIZOSPHERE. Chemicals produced and secreted by microorganisms that inhibit other microorganisms are called antibiotics. In principle, this is a form of allelopathy. Vascular plants react to microbial infection as well as to abiotic elicitors by synthesizing phytoalexins, which are low-molecular-weight, broad-spectrum antimicrobial compounds, another form of allelopathy. In contrast to phytoalexins, allelochemicals are produced without the action of an elicitor. See ANTIBIOTIC; PHYTOALEXINS; SOIL MICROBIOLOGY; TOXIN. Allelopathy is clearly distinguished from competition: in allelopathy a chemical is introduced by the plant into the environment, whereas in competition the plant removes or reduces such environmental components as minerals, water, space, gas exchange, and light. In the field, both allelopathy and competition usually act simultaneously, and assessing their individual influence and function is a continuing challenge. In controlled laboratory experiments with extracts or exudates from plants suspected of allelopathy, bioassays often yield significant allelopathic effects that are not manifest in nature. Soil and weather conditions, and the influence of surrounding plants and microbial interactions, carefully controlled in the laboratory, may mask allelopathic effects in the field. The unequivocal demonstration of allelopathy in the field as a result of the action of chemical compounds showing allelopathic activity in a laboratory bioassay is difficult because of the complexity and dynamic nature of the environmental interactions. Allelopathy was first observed in ancient times: Pliny (A.D. 23–79) related the inhibition of plant growth around walnut trees. The black walnut tree (Juglans nigra) produces a potent allelochemical, juglone (5-hydroxynaphthoquinone; structure 1); that inhibits growth of many annual O
OH
O
(1)
OH
OH
O-glucoside
(2)
469
470
Allelopathy herbs. Tomato and alfalfa under or near black walnut trees wilt and die. Alder (Alnus glutinosa), planted as nurse trees in black walnut orchards, decline and die from allelopathy within 8–13 years. Water leachate from leaves and extracts from the soil under a walnut tree contain juglone. Adding juglone to plant beds inhibits or kills the plants. Juglone is not, however, a general toxin for all plants; Kentucky bluegrass (Poa pratensis) and brambles such as wild raspberries and blackberries (Rubus spp.) grow under and around black walnut trees. In the cells of walnut trees, juglone is present in a safe form as 1,4,5-trihydroxynaphthalene glucoside (2), which is not inhibitory. Upon leaching from the leaves and bark by rain or release from the roots, the glucoside is hydrolyzed and the free hydroxynaphthalene is oxidized to the active allelochemical, juglone. The release of inhibitors, for example, hydrocyanic acid (HCN), from harmless precursors, such as cyanogenic glycosides, occurs frequently in allelopathy. Desert plants. For plants growing in habitats with extreme climates, such as a desert, competition for the limited resources is critical, and allelopathy may have survival value. Desert shrubs are often surrounded by a bare zone; thus, all the moisture of that zone remains available to the shrub and is not shared with other plants. In the Mojave Desert of California incienso (Encelia farinosa) inhibits the growth of desert annuals, which will grow only around dead incienso shrubs. From the decomposing leaf litter, 5acetyl-1-2-methoxybenzaldehyde (3) is released and O
hibition. While grazing by small mammals, birds, and insects—all taking advantage of the shrub for shelter—helps to keep down growth, grazing does not initiate or usually maintain the zones of inhibition. See HERBIVORY. The uphill and downhill inhibition zones of Salvia shrubs are about equal in size, indicating that watersoluble inhibitors are unlikely in the dry climate of the chaparral. This suggests the effect of volatile compounds. Volatile terpenes from Salvia and Artemisia and other shrubs and their surrounding soil were identified and found to be inhibitory to seed germination and plant growth. Camphor (4) and 1,8-cineole (5) are the most effective inhibitors; α-pinene (6), β-pinene (7), camphene (8), and thujone (9) are CH3 CH3 O CH3
CH3 CH3
O CH3
(5)
(4)
CH3
CH2
H3C
H3C CH3
CH3
(6)
(7)
C H
CH3O
CH3 O
COCH3
CH3
(3) persists in the desert soil, functioning as an allelochemical. Incienso is apparently not affected by its own toxin. See DESERT. Chaparral. Important allelochemicals in chaparrals are terpenoids and water-soluble compounds, such as organic acids and phenols. Terpenoids. The obvious and striking vegetation pattern of bare ground and stunted herbs and grasses around the shrub thickets of sagebrush (Salvia leucophylla) and wormwood (Artemisia californica) in the Californian chaparral is caused by an allelopathic effect of the volatile terpenes produced by these shrubs. Salvia and Artemesia are surrounded by a zone of bare soil 3–6 ft wide (1–2 m), followed, toward the periphery, by a zone of annual herbs and grasses with stunted growth and limited development and, finally, thriving grassland with wild oats (Avena fatua), brome grass (Bromus rigidus), and other grass species. Careful analyses have eliminated shade, nutrient deficiency, water availability, root density, slope of ground, and animal feeding and running activity as the causes for growth in-
CH3 CH2
(8)
CH3 CH3
(9)
also active volatile allelochemicals. Growth inhibition is most obvious around shrubs growing in the fine-textured Zaca clay, which absorbs the terpenes and thus facilitates contact with the roots. Spectacular fires, which sweep the chaparral approximately every 25 years, volatilize the soil-bound terpenes and destroy the shrubs. The burned sites quickly turn green with annual herbs and grasses that dominate for approximately 6 years until the shrubs grow back from surviving underground parts. With the accumulation of terpenes in the soil, the zones of bare ground and stunted growth reappear and remain apparent until the next fire. Terpenoids are common plant products and might be expected to serve as allelochemicals in other plant species. Indeed, for eucalyptus trees and sassafras, which do not inhabit the chaparral, terpenoids have been identified as effective allelochemicals; in other
Allelopathy plants, especially conifers, terpenoids are suspected of acting as such also. See TERPENE. Water-soluble allelochemicals. Chamise (Adenostema fasciculatum) and manzanita (Arctostaphylos glandulosa), two common chaparral shrubs species, lack volatile terpenes, but because growth around them is also inhibited, other allelochemicals may be involved. Ferulic acid (10), p-hydroxybenzoic acid (11), and vanillic acid (12) are found in their leaves
HO HO
COOH HO
(15) O HO
HO
CH
CH
C H
COOH
(16) H3CO
(10)
CH3
CH
OH
C
C
C
C
CH
OH
COOCH3
(17) CH3
HO
HO
COOH
COOH
H3CO
(11)
(12)
as well as in the surrounding soil and act as strong inhibitors of seed germination and growth. Other phenols and phenolic acids are present in the leaves but are not always detected in the soil. Rainfall in the chaparral is only moderate. The coastal fog, which prevails throughout the year, however, supplies sufficient moisture to leach the phenolic compounds into the soil, rendering them effective allelochemicals. See CHAPARRAL. Old-field succession. Allelopathy has been studied in the four stages of old-field succession of the Oklahoma-Kansas prairie. The pioneer weed stage, lasting only 2–3 years, is characterized by sturdy weeds with low mineral needs, such as sunflower (Helianthus annuus), ragweed (Ambrosia psilostachya), Johnson grass (Sorghum halepense), crab grass (Digitaria sanguinalis), and spurge (Euphorbia supina). Autoallelopathy of the pioneer weeds accelerates their own replacement by the dominant species of the second stage: namely the triple-awn grass (Aristida oligantha), which lasts about 9–13 years. Tolerant of the pioneer weed-produced allelochemicals, such as chlorogenic acid (13), p-coumaric acid (14), gallic acid (15), and p-hydroxybenzaldehyde (16), e-awn grass also HO
HO
O
HO
CH
CH
C
OH
O COOH HO
(13)
HO
CH
(14)
CH
COOH
CH2
CH2
C
C
C
C
CH
CH
COOCH3
(18) has low mineral requirements. Allelochemicals produced from triple-awn grass, such as gallic acid and tannic acid, inhibit free-living and symbiotic nitrogen-fixing bacteria, thus giving the grass a selective advantage over invading species of the third stage that require higher nitrogen concentrations. Inhibition of the nitrifying bacteria by allelochemicals produced by the invading species, in turn, controls the rather slow succession of the third and fourth stages. Gradually, ammonia in the soil increases to such a level that species of the third stage, perennial bunchgrasses—lasting some 30 years—and finally, the climax prairie, can slowly claim the field. In Japanese urban wasteland, weed succession, which is similar to that of old fields, appears to be controlled by polyacetylene-type allelochemicals. Matricaria esters (17) and lachnophyllum ester (18) in the pioneer weeds, Solidago altissima and Erigeron species, have also been identified in the soil, suggesting ecological significance. See ECOLOGICAL SUCCESSION. Aqueous habitats. Allelopathy is not restricted to terrestrial habitats but occurs also in fresh water and marine environments. The spikerush (Eleocharis coloradoensis) is an aquatic plant that has been shown to produce an allelochemical effect. Allelopathic substances in saltwater are of special interest with regard to the red tide, that is, blooms of microorganisms that discolor the water and discharge toxins. Algae-produced allelochemicals may be able to limit the growth of red tide producers. Nannochloris, an alga, synthesizes a cytolytic agent affecting a red tide–causing dinoflagellate. Other algae have been found to secrete fatty acids that inhibit algal growth. Agriculture. The reduction of crop yields by weed competition is aggravated by the allelopathic effect of weeds on crops. It was already noted in 1832 that thistles (Cirsium) diminished the yield of oats, and spurge (Euphorbia) diminished the yield of flax. The growth inhibition of wheat by couch grass (Agropyron repens) results from impaired phosphate uptake,
471
472
Allelopathy even in the presence of adequate phosphate supply, and results in severe crop losses. Leachates from Italian ryegrass (Lolium multiflorum), its litter, and the soil in which it grows inhibit the growth of oat, clover, lettuce, and bromegrass. Foxtail species (Setaria) reduce the growth of corn by as much as 35% through allelopathy. The combined effects of foxtail allelopathy and competition can result in 90% growth reduction. Crop plants may inhibit their own growth and reduce the yield of subsequent crops, an observation well known to farmers as soil sickness long before the investigations of allelopathy. Crop rotation cures soil sickness only when the subsequent crop is not affected by the accumulated allelochemicals of the previous crop or when they have been detoxified by soil microorganisms. Cabbage, rich in thiocyanates, inhibits grass species, tobacco, and grapes. Woad (lsatis tinctoria) contains high levels of glucosinolates (mustard oil glycosides), especially glucobrassicin (19), and its planting is thereN CH2
OSO3
C S-glucosyl
N
(19)
fore often avoided because of its ill effects on succeeding crops. Straw from wheat, corn, rye, rice, or sorghum decomposing in the field yields phenolic compounds, such as vanillic acid (12), ferulic acid (10), p-hydroxybenzoic acid (14), p-coumaric acid (13), o-hydroxyphenylacetic acid, and protocatechuic acid (20), all of which exert allelopathic ef-
HO
COOH HO
(20)
fects. In addition, an increase in the soil microflora in response to the crop residue generally results in nitrogen immobilization and, hence, in reduced growth of the succeeding crop. As usual in nature, competition and allelopathy occur simultaneously, and the relative effect of each on plant growth and yield is difficult to assess. Soil microorganisms have a critical role in the allelopathy responsible for clover soil sickness. Clover roots exude toxic isoflavonoids, which do not cause allelopathy but are converted by microorganisms to salicylic acid (21), p−methoxybenzoic acid (22), 2,4hydroxybenzoic acid (23), and other phenolic com-
COOH OH
(21)
OH H3CO
COOH
HO
(22)
COOH
(23)
pounds that are responsible for clover soil sickness. Soil microorganisms may also prevent allelopathy. In sandy soils, sorghum autoallelochemicals reduce the yield of the first and subsequent sorghum crops. In clay-rich soils, sorghum autoallelopathy does not become manifest because fungi degrade the allelochemicals before growth of the new seedlings starts. In the complex interaction between crops and weeds and among crops lie exciting possibilities for utilizing the allelopathic properties of crop plants for weed control. Allelopathic weed control with crop plants has long been practiced by farmers and horticulturists. The challenge is twofold: to minimize the negative impact of allelochemicals on crop growth and yield and to exploit allelopathic mechanisms for pest control and crop growth regulation strategies. Allelochemicals in new crop cultivars may provide naturally occurring pesticides that can limit or suppress weeds as well as prevent insect and nematode attack and damage. The biotechnology resources for the production of herbicide-resistant crops could then be channeled into the engineering of other desired crop qualities. Allelochemicals may furnish an entirely new generation of naturally produced weedcontrolling compounds, replacing synthetic herbicides and pesticides with nonaccumulating, easily degradable substances. See HERBICIDE. The natural occurrence of many plant species in pure stands that apparently cannot be invaded by other species may result, in part, from allelopathic mechanisms that could be exploited to achieve weed-free crop stands. Barley, considered a smother crop, was thought to inhibit weeds through competition with its extensive root system. However, even in the absence of competition, barley inhibits weeds, such as chickweed (Stellaria media), shepherd’s purse (Capsella bursa-pastoris), and tobacco, but not wheat. Gramine (24), an alkaloid, which occurs CH3 N CH3
N H
(24) in barley and inhibits Stellaria growth, may function as an allelochemical. Allelochemicals in seeds, including alkaloids (for instance, in coffee beans), may protect seeds from microbial attack and destruction while they remain viable and buried in the soil for years. Agricultural practice uses cover crops, which are either killed by frost or by desiccant sprays, and the allelochemicals of the crop residues suppress weed growth. Oats, corn, sunflower, sorghum, and fescue appear promising in allelopathic weed control. See ALKALOID.
Allelopathy Not all weed-crop interactions, however, are necessarily detrimental for the crop. Corn cockle (Agrostemma githago), growing in wheat fields in former Yugoslavia, increases the wheat yield but suppresses its own growth. Agrostemin and gibberellin have been identified as the allelochemicals of corn cockle. Triacontanol, CH3(CH2)28CH2OH, an allelochemical from alfalfa, promotes growth of tomato, cucumber, and lettuce among other species. See GIBBERELLIN. Parasitism. Infestation of vascular plants by vascular plant parasites, often treated as an oddity, is an allelopathic relationship. The seeds of the parasite witchweed (Striga asiatica) germinate in response to root exudate from host plants. While the active compounds have not yet been identified, strigol (25) induces witchweed seed germinaH3C
CH3
O
O
CH
CH3 O
O
O OH
(25) tion and the formation of a haustorium (a nutrientabsorbing structure) that attaches to the host. Strigol is beneficial to the parasite by stimulating its growth and facilitating recognition of the host, but is detrimental to the plant producing it, when the parasite finds the producer and attaches itself. Purple false foxglove (Agalinis purpurea), which parasitizes legumes, responds to xenognosin A (26) with haustorium formation. HO
OCH3
OH
(26) Forestry. As in agriculture and horticulture, allelopathy along with other ecologically important factors affects the pattern of forest growth and development and therefore has strong economic implications. While examples of tree allelopathy, such as that of the black walnut, were well known, the general significance of allelopathy for forestry has been recognized only in recent times. Oak trees left as seed trees may impede reforestation by retarding the growth of the tree seedlings. Salicyclic acid (21) has been identified as one of the allelochemicals, but strong competition from oak roots and physical suppression of seed germination by oak litter also take part in retarding growth. Bare areas under hackberry (Celtis laevigata) result from accumulation of the same phenolic allelochemicals that affect many other plant communities. Synergistic action of these allelochemicals, while often suspected, has been observed with hackberry toxins. Seasonality of allelopathy in a forest of sycamore (Platanus occidentalis), rough-leaved hackberry (C. occidentalis), and white oak (Quercus
alba) may indicate that certain allelochemicals are effective upon release, whereas others may require chemical change by microorganisms before they become active. Production and release of toxins must exceed decomposition for allelopathy to occur. Grevillia robusta, a stately tree in the subtropical rainforest of Australia, does not form pure stands, and, not surprisingly, does not grow in monoculture beyond 10–12 years, when it declines because of strong autoallelopathy. Bracken stands or bracken-grassland communities surrounded by forest often remain free of trees. Phenolic compounds, including caffeic acid (27) and HO HO
CH
CH
COOH
(27) ferulic acid (10), appear to be the controlling allelochemicals. A dense ground cover of bracken, wild oat grass (Danthonia compressa), goldenrod (Solidago rugosa), and flat-topped aster (Aster umbellatus), does not allow a forest in the Allegheny Plateau of northwestern Pennsylvania to grow back, and half a century after clear-cutting, only a few scattered black cherry (Prunus serotina) and red maple (Acer rubrum) grow. Other woodland ferns, such as the hay-scented fern (Dennstaedtia punctilobula) and New York fern (Thelypteris noveboracensis), also contribute to the allelopathic effect. Release of allelochemicals from bracken appears to be timed for the most effective suppression of surrounding vegetation. In the tropics, toxins are released from green bracken leaves year-round, while in southern California and in the Pacific Northwest, toxins are released from dead bracken fronds, litter, and roots when seeds germinate, either at the beginning of the rainy season or in the spring, respectively. Allelopathic inhibition of forest regrowth may be indirect by affecting the mycorrhizal fungi (fungi which grow mutualistically in association with roots of vascular plants). Shrub litter, heather (Calluna vulgaris), and prairie grasses inhibit Norway spruce (Picea abies) development by preventing the establishment of mycorrhiza. See FOREST AND FORESTRY; MYCORRHIZAE. Action mechanisms of allelochemicals. While numerous allelochemicals have been identified, most studies on allelopathy deal with unidentified compounds. Allelochemicals are secondary plant products and belong to a few major classes of chemical compounds: terpenoids, phenolic compounds, phenylpropane derivatives, flavonoids, long-chain fatty acids, organic cyanides, probably also alkaloids and purines, and perhaps some steroids. In the field, allelopathy involves a complex of compounds; a single specific compound does not appear to produce the observed allelochemical effects on neighboring plants. This complicates the investigation of the action mechanism of allelochemicals.
473
474
Allergy There is no single set of physiological reactions controlled by allelochemicals that result in allelopathy; the action of allelochemicals is diverse and affects a large number of physiological functions and biochemical reactions. At the whole-plant level, allelochemicals interfere with internode elongation, leaf expansion, cell division, dry-weight accumulation, as well as seed germination. Growth responses mediated by natural plant growth regulators, such as auxin or gibberellin, are often impaired. In field experiments, allelochemicals from some crop residues promote bacterial and fungal lesions on roots, enhancing diseases of the test seedlings. See PLANT HORMONES. Common to many allelochemicals is their effect on membrane permeability. Ion uptake can be enhanced or reduced; certain phenolic acids inhibit the uptake of phosphate and potassium. Fescue leachates impair uptake of phosphorus and nitrogen in sweetgum. Some flavonoids, phenolic acids, and juglone inhibit membrane adenosine triphosphatase activity that is involved in the control of membrane transport. Sphagnum acidifies its habitat by releasing hydrogen ions, which may make conditions less suitable for competing plants. Significant effects on photosynthesis by allelochemicals have been observed. Allelochemicals from pigweed and foxtail litter affect photosynthesis in corn and soybeans and, in addition, alter the biomass partitioning into the leaves compared with that of the total plant. Stomatal movement, which regulates gas exchange and water status, and net photosynthesis may also be inhibited. At the cellular and subcellular level, kaempferol (28), a flavonol, interferes with the OH O
HO
OH OH
O
(28) photosynthesis of isolated chloroplasts by inhibiting coupled electron transport and photophosphorylation (adenosine triphosphate synthesis). In isolated mitochondria, kaempferol inhibits oxidative phosphorylation. Respiration was also inhibited by juglone (1) and other allelochemicals. The biosynthesis and metabolism of proteins, lipids, and organic acids are also affected by allelochemicals. While it is relatively easy and straightforward to isolate and identify allelochemicals, to test them in bioassays, and to measure their physiological and biochemical reactions, in the field it is difficult to prove directly that the same compounds move from the producing plant to the neighboring plants and function as allelochemicals in the field. Manfred Ruddat Bibliography. W. Y. Garner and J. Harvey, Jr. (eds.), Chemical and Biological Control in Forestry, ACS Symp. Ser. 238, 1984; M. B. Greene and P. A. Hedin (eds.), Natural Resistance of Plants to Pests: Role of Allelochemicals, ACS Symp. Ser. 296, 1986; J. B.
Harbone, Introduction to Ecological Biochemistry, 4th ed., 1994; M. M. Harlin, Allelochemistry in marine macroalgae, CRC Crit. Rev. Plant Sci., 5:237– 249, 1987; A. R. Putnam and Chung-Shih Tang (eds.), The Science of Allelopathy, 1986; E. L. Rice, Allelopathy, 2d ed. 1984; A. C. Thompson (ed.), The Chemistry of Allelopathy: Biochemical Interactions among Plants, ACS Symp. Ser. 268, 1985; G. R. Waller (ed.), Allelochemicals: Role in Agriculture and Forestry, ACS Symp. Ser. 330, 1987.
Allergy Altered reactivity in humans and animals to allergens (substances foreign to the body that cause allergy) induced by exposure through injection, inhalation, ingestion, or skin contact. The most common clinical manifestations of allergy are hay fever, asthma, hives, atopic (endogenous) eczema, and eczematous skin lesions caused by direct contact with allergens such as poison ivy or certain chemicals. Allergens. A large variety of substances may cause allergies: pollens, animal proteins, molds, foods, insect venoms, foreign serum proteins, industrial chemicals, and drugs. Most natural allergens are proteins or polysaccharides of moderate molecular size (molecular weights of 10,000 to 200,000). Chemicals or drugs of lower molecular weight (haptens) have first to bind to the body’s own proteins (carriers) in order to become fully effective allergens. For the development of the hypersensitivity state underlying clinical allergies, repeated contact with the allergen is required. Duration of the sensitization period is usually dependent upon the sensitizing strength of the allergen and the intensity of exposure. Some allergens (for example, saliva, urine, and hair proteins of domestic animals) are more sensitizing than others. In most instances, repeated contact with minute amounts of allergen is required: several annual seasonal exposures to grass pollens or ragweed pollen usually occur before an overt manifestation of hay fever. On the other hand, allergy to cow milk proteins in infants can develop within a few weeks. When previous contacts with allergens have not been apparent (for example, antibiotics in food), an allergy may become clinically manifest even upon the first conscious encounter with the offending substance. Besides the intrinsic sensitizing properties of allergens, individual predisposition of the allergic person to become sensitized also plays an important role. Clinical manifestations, such as hay fever, allergic asthma, and atopic (endogenous) dermatitis, occur more frequently in some families: the inheritance of a capacity to develop these forms of allergy has been called atopy or the atopic state. These terms indicating a disease “not at its proper place” are misnomers but are sanctioned by historical usage. In other clinical forms of allergy, genetic predisposition, though possibly present as well, is not as evident. Mechanisms. Exposure to sensitizing allergens may induce several types of immune response,
Allergy and the diversity of immunological mechanisms involved is responsible for the various clinical forms of allergic reactions which are encountered in practice. Three principal types of immune responses are encountered: the production of IgE antibodies, IgG or IgM antibodies, and sensitized lymphocytes. See ANTIBODY; IMMUNOGLOBULIN. IgE antibodies. IgE antibodies belong to a peculiar class of serum immunoglobulins which is responsible for the majority of allergies of the so-called immediate or anaphylactic type. IgE antibodies are present in blood in very small amounts (normally less than 0.2 microgram per milliliter), but have the capacity to bind strongly to the membrane of tissue mast cells and blood basophils. These cells are the major effector components in an immediate allergic reaction. They contain a wide variety of preformed inflammatory mediators, such as histamine and serotonin. They also have the capacity to form, upon interaction with allergen, further substances causing inflammatory changes in tissue. The latter include increased vascular permeability that results in swelling and is manifested by redness and wheals as observed in hives. Among such mediators formed after cell stimulation, leukotrienes, the platelet-activating factor (PAF), and several other substances which attract other blood cells, such as eosinophils, to the site of reaction have gained the most attention. Mast cells and basophils may be compared to powder kegs which will be fired upon the interaction of allergen with its corresponding cell surface–bound IgE antibody. Understandably, such a reaction proceeds quite rapidly in a sensitized individual following renewed contact with allergen. Within minutes, hay fever patients develop symptoms after inhaling grass pollen, as do individuals allergic to bee venom after being stung. Generally, individuals possessing a genetic predisposition to atopic diseases are extremely prone to develop IgE antibodies to a variety of inhaled or ingested allergens. The IgE antibodies may also bind, although with lesser affinity, to some mononuclear cells and in particular to Langerhans cells in the skin. This phenomenon may be involved in the mechanism of some eczematous skin lesions, as seen in the frequent allergic condition termed atopic dermatitis. See HISTAMINE; SEROTONIN. IgG or IgM antibodies. The production of IgG of IgM antibodies requires, as a rule, more massive exposure to allergens than that leading to the formation of IgE antibodies. Upon renewed contact with sufficient amounts of allergen, immune complexes may form with serum antibodies and be deposited in various tissues where they cause an acute inflammatory reaction. Preferential sites for such reactions, depending upon the mode of entry and distribution of the allergen, are alveolar walls of the lungs, glomeruli of kidneys, synovial membranes of joints, and walls of small blood vessels of the skin. Immune complex deposition below the basement membrane of blood vessels induces a cascade activation of complement (a multicomponent system of interacting blood proteins) which generates further active mediators of inflammation. The most important include
the fragments of complement factors C3 and C5 (denominated C3a and C5a or anaphylatoxins) which are chemotactic (attracting) for white blood cells, especially polymorphonuclear leukocytes, attracting them to infiltrate the site of the reaction. Once arrived there, these cells proceed to destroy excess allergen and immune complexes, and by various mechanisms enhance the inflammatory reaction. The process takes several hours, and reactions of this type (Arthus reactions) become clinically apparent and peak only 6 to 12 h following contact with allergen. See IMMUNE COMPLEX DISEASE; INFLAMMATION. Sensitized lymphocytes. Sensitized lymphocytes are responsible for the clinical manifestations of so-called delayed hypersensitivity. If an individual has been exposed to allergen and has developed sensitized lymphocytes capable of reacting with it, renewed exposure will trigger the lymphocytes to produce an array of glycoproteins called lymphokines. Lymphokines are also mediators capable of causing considerable inflammation in the surrounding tissues. Although produced at first only by the relatively small proportion of lymphocytes which specifically recognize and interact with allergen, lymphokines have amplifying effects. They stimulate, indiscriminately, other cell types, such as the macrophages and monocytes, attract these cells to the site of the reaction, and cause the formation of an inflammatory infiltrate. In the skin, such a reaction manifests itself by redness, swelling, and the formation of an inflammatory papule. A typical example is the inflammatory nodular reaction experienced upon injection in the skin of protein allergens produced by tubercle bacilli (tuberculin reaction). Such a reaction occurs only in people who were once infected with tubercle bacilli, and who have developed an immune response to these microorganisms. Skin reactions occurring upon contact withbreak sensitizing agents such as poison ivy or industrial chemicals also are caused by the same immune mechanism. Since production of lymphokines by lymphocytes following stimulation by allergens takes several hours and since the reaction becomes clinically apparent only when a minimal number of infiltrating cells (mostly monocytes) have accumulated at the reaction site, the time elapsed between allergen contact and clinical reaction is usually 24 to 72 h. Such reactions are therefore termed delayed reactions, in contrast to the immediate (anaphylactic) reactions which occur within minutes after contact with allergen C. See HYPERSENSITIVITY. Clinical forms. Among the clinical manifestations of allergy are allergic rhinitis, bronchial asthma, food allergy, occupational allergy, skin allergy, and anaphylactic shock. Allergic rhinitis. The most frequent manifestation of allergy is undoubtedly allergic rhinitis. Allergic rhinitis was originally misnamed hay fever because it was believed that hay causes the disorder, and “fever” was a term loosely applied to many ailments, even when not accompanied by fever. However, in addition to nasal symptoms (watery nasal discharge, sneezing, runny eyes, and itchy nose, throat, and roof of the
475
476
Allergy mouth), the individual may experience fatigue, irritability, and loss of appetite. In industrialized countries, the frequency of allergic rhinitis has been estimated to be 5 to 8% of the population and appears to be increasing. Allergic rhinitis may be seasonal, as the traditional hay fever, or perennial, occurring year-round. The seasonal form is usually caused by pollens from grasses, weeds, or trees. Related symptoms appear only during that time of the year when the pollens to which the individual is sensitive occur in the air. When allergens are present all the time, allergic rhinitis may occur year-round. Frequent causes are allergens present in the house or workplace, such as dust, molds, and animal danders. Chronic or recurrent rhinitis, however, is not always related to allergy; in many instances it may be caused by an abnormal irritability of the nose (socalled intrinsic or vasomotor rhinitis) attributable to still ill-defined factors, such as hormonal factors or stress. Bronchial asthma. The allergic disease which by its frequency and severity causes the most problems is bronchial asthma. Although an attack of asthma may be easily described, since it involves acutely occurring shortness of breath and wheezing, a general accurate definition, encompassing all possible causes and mechanisms of the disease, is still not available. Constriction of tiny bronchi in the lung airways, accompanied by an augmentation of bronchial secretions and the swelling of bronchial mucosa, contributes to obstruct the free passage of air through the lungs, producing wheezing. However, wheezing is not always indicative of asthma. Furthermore, in numerous instances, the cause of asthma is not related to allergy against some airborne allergens. It is customary to classify asthma as extrinsic, that is, due to contact with an external allergen, or intrinsic, that is, due to endogenous causes which are still poorly defined, such as stress or bronchial hyperreactivity owing to hormonal and other factors. However, even if the causes of disease may be very dissimilar, the mechanisms and particularly the mediators responsible for allergic inflammation in the lung may be the same, leading to almost indistinguishable clinical pictures. The same airborne allergens which cause allergic rhinitis may also be responsible for seasonal or yearround asthma. Indeed, allergic rhinitis and asthma are often associated in the same individual. In addition, ingestion of allergens in foods and inhalation of drugs by aerosol or of sensitizing industrial chemicals at the workplace are frequent causes of allergic asthma. Allergic asthma is considered a major public health problem owing to its severity, and to the possibility of death during an acute and prolonged asthma attack (status asthmaticus), and also because of the possible evolution to severe chronic respiratory insufficiency with invalidism in neglected cases. Available statistics indicate that in most countries 2–4% of the population is affected by asthma. In view of the
large variety of possible causes of asthma, different mechanisms may be operating. In classical cases of allergic asthma due to airborne allergens, the presence of IgE antibodies probably plays a major role, but it is not excluded that other forms of immune reactions also occur. See ASTHMA. Food allergy. This may occur as a reaction to natural products (for example, celery, milk, and egg proteins) as well as to chemicals or preservatives added to foods (for example, dyestuffs and antioxidants). Allergy to foods may manifest itself at the primary site of contact with allergen, namely in the gastrointestinal tract, causing symptoms such as diarrhea, abdominal pain, nausea, and vomiting. However, it may also take more general forms once the allergen has been absorbed, causing symptoms in the upper airways (rhinitis and asthma) and in the skin (hives), and, in cases of extreme hypersensitivity, generalized anaphylactic shock. In infants from allergic families, the development of allergy to cow milk is quite frequent; strict breast feeding during the first months of life is therefore increasingly advocated. The breast-feeding mother should abstain from ingesting large quantities of potential allergens such as milk and eggs, since the food allergens can pass into the mother’s milk and sensitize the child. After a few months, the susceptibility of infants to sensitization by ingestion appears to decrease markedly. Occupational allergies. Such allergies become increasingly frequent in industrialized countries. Chemicals used in plastic and pharmaceutical industries or in metal refining and construction work are frequent offenders. However, natural organic materials, such as baker’s flour, animal products, castor bean, silk, and furs, may be involved as well. Here, too, the most frequent manifestations involve the upper respiratory tract and the skin. A peculiar form of occupational lung disease, hypersensitivity pneumonitis, is characterized by repeated episodes of cough, fever, chills, and shortness of breath, occurring 4–10 h after massive exposure to some airborne allergen. In contrast to allergic asthma, the lesions do not occur at the level of the bronchi but of the alveolar walls in the lung. The inflammatory reaction taking place there may have severe consequences when occurring repeatedly, since it leads to progressive fibrosis of the lung, impairment of gas exchanges with blood, and respiratory insufficiency. The mechanisms of this type of allergic reaction are not entirely elucidated, but it is known that IgE antibodies are not involved. IgG antibodies are frequently present and possibly cause the formation of immune complexes with allergen. However, delayed hypersensitivity mechanisms are probably at play as well. Among the most frequent offending allergens are molds sometimes found in hay (causing “farmer’s lung”) and in various other locations, such as air conditioners, cellars, and logs used by woodworkers. Skin allergy. This condition may take several forms. Atopic dermatitis is a chronic superficial itching
Allergy inflammation of the skin (eczema) which occurs in 0.5–0.7% of the population and is often of long duration. Atopic dermatitis usually occurs early in life together with cow milk allergy (infantile eczema). Its evolution is unpredictable, and the disease may seriously affect psychological health and working ability of the individual. IgE antibodies frequently are notably increased. This disease exhibits the same genetic predisposition (“atopy”) as allergic rhinitis and asthma. Allergic contact dermatitis is characterized by a red rash, swelling, intense itching, and sometimes blisters at the site of contact with the allergen, which may be of plant origin (for example, poison ivy or turpentine) or may be an industrial chemical. Lesions are caused by sensitized lymphocytes and belong to the delayed hypersensitivity type of allergy. Urticaria (hives) is a very common manifestation of skin allergy. When evoked by sensitization to some allergen, IgE antibodies are present. Foods, drugs such as the penicillins, and, more exceptionally, natural inhalation allergens such as pollens are the most frequently involved. However, utricaria, resulting from the release of inflammatory mediators (particularly histamine) by skin mast cells, is even more frequently caused by nonallergic mechanisms. Some foods (for example, crustaceans or strawberries) contain substances which, when ingested in sufficient quantities, can directly trigger mast cells without prior sensitization and formation of IgE antibodies. Anaphylaxis. The most feared allergic reaction is an acute event known as anaphylactic shock. This is an IgE-mediated reaction with sudden onset: symptoms develop within minutes after exposure to the triggering allergen. The victim may experience a feeling of great anxiety; symptoms include swelling and skin redness often accompanied by hives, vomiting, abdominal cramps and diarrhea, life-threatening breathing difficulties due to swelling of the throat, and severe bronchospasm. A drastic fall in blood pressure may cause fatal shock within a few minutes. Insect stings and drugs such as the penicillins are now the most frequent causes of anaphylactic shock; in earlier times the therapeutic injection of animal serum proteins (for example, horse antitetanus serum) was often responsible. Some acute reactions, called anaphylactoid reactions, may cause similar symptoms but are not due to IgE-mediated sensitization. Such reactions are encountered mostly in people intolerant to some drugs, such as aspirin, or to contrast media used in diagnostic radiology. See ANAPHYLAXIS. Diagnosis. Diagnosis of allergic diseases encompasses several facets. The large variety of clinical forms of allergic reactions and the multitude of allergens possibly involved frequently require prolonged observation and in-depth investigation of the individual’s life habits. Since many clinical manifestations of allergy are mimicked by nonallergic mechanisms, it is usually necessary to use additional diagnostic procedures to ascertain whether the person has developed an immune response toward the incriminated aller-
gen. Such procedures primarily consist of skin tests, in which a small amount of allergen is applied on or injected into the skin. If the individual is sensitized, a local immediate reaction ensues, taking the form of a wheal (for IgE-mediated reactions), or swelling and redness occur after several hours (for delayed hypersensitivity reactions). The blood may also be analyzed for IgE and IgG antibodies by serological assays, and sensitized lymphocytes are investigated by culturing them with the allergen. Since the discovery of the responsible allergens markedly influences therapy and facilitates prediction of the allergy’s outcome, it is important to achieve as precise a diagnosis as possible. Most of the tests described above indicate whether the individual is sensitized to a given allergen, but not whether the allergen is in fact still causing the disease. Since in most cases the hypersensitive state persists for many years, it may well happen that sensitization is detected for an allergen to which the individual is no longer exposed and which therefore no longer causes symptoms. In such cases, exposition tests, consisting of close observation of the individual after deliberate exposure to the putative allergen, may yield useful information. Therapy. Therapy of allergic reactions may be performed at several levels. Whenever feasible, the most efficient treatment, following identification of the offending allergen, remains elimination of allergen from the person’s environment and avoidance of further exposure. This treatment is essential for allergies caused by most household and workplace allergens. Inflammatory reactions and corresponding symptoms caused by allergic mediators may be relieved by several drugs acting as pharmacological antagonists to these mediators, for example, antihistaminics or steroids. Other drugs act at a more central level and influence the ease with which mast cells or basophils release their mediators; this is apparently the case for chromones. Since the state of hypersensitivity usually remains present over many years, especially when contact with allergen is periodically renewed (as is the case with pollens), such an approach to allergy therapy may be only palliative. Attempts to influence the hypersensitive state itself in IgE-mediated allergy are more decisive. Repeated injections of increasing doses of allergen progressively diminish the degree of hypersensitivity. This procedure, which must usually be pursued over several years, is referred to as hyposensitization or immunotherapy. See IMMUNOLOGY; IMMUNOTHERAPY. A. L. de Weck Bibliography. R. Davies and S. Ollier, Allergy: The Facts, 1989; H. F. Krause, Otolaryngic Allergy and Immunology, 1989; L. M. Lichtenstein and A. S. Fauci, Current Therapy in Allergy, Immunology, and Rheumatology, 5th ed., 1996; G. Melillo and J. A. Nadel (eds.), Respiratory Allergy, Advances in Clinical Immunology and Pulmonary Medicine, 1993; D. Reinhardt and E. Schmidt (eds.), Food Allergy, 1988; I. M. Roitt, Essential Immunology, 1991.
477
478
Alligator
Alligator A large aquatic reptile of the family Alligatoridae. Common usage generally restricts the name to the two species of the genus Alligator. The family also includes three genera (five species) of caiman. With the crocodiles and gharial (also spelled gavial), these are survivors of archosaurian stock and are considered close to the evolutionary line which gave rise to birds. Alligator species. The two living species have a disjunct subtropical and temperate zone distribution. The American alligator (A. mississippiensis) ranges throughout the southeastern United States from coastal North Carolina (historically from southeastern Virginia) to the Rio Grande in Texas, and north into southeastern Oklahoma and southern Arkansas (see illustration). Poaching and unregulated hunting for the valuable hide decimated the alligator populations until the animal was placed on the U.S. Endangered Species List in 1967. The species has responded well to protection and has become abundant in many areas in its range, particularly in Florida and parts of Georgia, Louisiana, and Texas, where it is now a common sight in the freshwater habitats, including swamps, marshes, lakes, rivers, and even roadside ditches. The second species is the Chinese alligator (A. sinensis), restricted to the region of the Yangtze River valley, where it inhabits burrows in the floodplains and riverbanks. This species is considered “critically endangered” by international conservation organizations. It is believed that less than 200 individuals remain in the wild. The American alligator is by far the larger of the two species, exceeding 15 ft (4.5 m). The average length of A. sinensis is 4–5 ft (1.2–1.5 m). Caiman species. There are five species of caimans, once often sold as “baby alligators.” Caimans differ from alligators in technical details of their internal anatomy and scale characteristics. The black caiman (Melanosuchus niger) of the Amazon Basin of South America strongly resembles the American alligator in superficial appearance and may reach a length approaching 16 ft (5 m). The spectacled caiman (Caiman crocodilus), with its several subspecies, is sold in the pet trade. It is the
Half-submerged American alligator (Alligator mississippiensis). (Photo courtesy of W. Ben Cash)
most widely distributed of the caimans, ranging from southwestern Mexico to Paraguay and Bolivia. Spectacled caimans reach a length of 8 ft (2.5 m) and can be distinguished from alligators by the presence of a ridge across the snout between their eyes, the nose bridge of the “spectacles.” The broad-nosed caiman (C. latirostris) occurs from southern Brazil to Paraguay and northern Argentina and reaches a length of 9 ft (3 m). The two species of smooth-fronted caiman, Paleosuchus trigonatus and P. palpebrosus, are the smallest of the family, reaching a maximum length of 7 ft (2 m). They occur over most of northern South America, south to Peru, Bolivia, and southeastern Brazil. Anatomical features. Alligators, including caimans, are generally distinguished from crocodiles by their broader, rounded, and more massive snout, the arrangement of their teeth, and other technical anatomical details. As in all crocodilians, the teeth are conical and sharp, equivalent in shape (homodont), and replaced throughout life. Moreover, as in all crocodilians, alligators possess a functionally fourchambered heart, a secondary palate, a septum or “diaphragm” that separates the lung and peritoneal regions, a compressed tail, webbed feet, and other adaptations for an aquatic existence. Nesting. All alligators build a conical nest of available materials near the water’s edge, where they deposit 15–60 leathery-shelled eggs. Incubation takes 30–60 days; females of at least some species may guard the nest, assist in the hatching process, and remain with the young for as long as 3 years. In some cases, females carry young in their mouth to nearby water. Habitat and behaviors. Alligators are generalized carnivores, feeding on any invertebrates, fishes, reptiles and amphibians, birds, or mammals that can be caught and overpowered. They play a dominant role in the energy and nutrient cycling in their habitats, and through their construction of “gator holes,” which often retain water in times of drought, are considered important to the survival of many others species. Alligators have also been found to have a complex social system that relies heavily on communication. Visual, tactile, auditory, and even subauditory communications have been described in alligator populations. Males roar during the breeding season to attract females, and females bellow and growl to communicate receptiveness to males. Using trunk muscles and low-frequency sound, alligators create rapid vibrations on the water surface to communicate. The vibrations coupled with the low-frequency sound can be heard at great distances (perhaps up to a kilometer). Males use a diversity of sounds as a threat warning, including a “cough” and hiss. Juveniles employ distress calls when they feel threatened. These calls attract the attention of the nearby female. See ARCHOSAURIA; CARDIOVASCULAR SYSTEM; CROCODILE; CROCODYLIA; DENTITION; REPTILIA. W. Ben Cash; Howard W. Campbell Bibliography. S. Grenard and W. Loutsenhizer, Handbook of Alligators and Crocodiles, Krieger,
Allometry 10,000 1000 100 organ mass, kg
Melbourne, FL, 1995; S. B. Hedges and L. L. Poling, A molecular phylogeny of reptiles, Science, 283:998–1001, 1999; D. W. Linzey, Vertebrate Biology, McGraw-Hill, New York, 2001; F. H. Pough et al., Herpetology, 3d ed., Prentice Hall, Upper Saddle River, NJ, 2004; K. A. Vliet, Social dynamics of the American alligator (Alligator mississippiensis), Amer. Zool., 29:1019–1031, 1989; G. R. Zug, L. J. Vitt, and J. P. Caldwell, Herpetology, Academic Press, San Diego, 2001.
10
skeleton = 0.06m 1.08 lung = 0.01m 0.99
1 0.1 0.01
brain = 0.01m 0.70
0.001
heart = 0.58m 0.98
0.0001 0.001 0.01 0.1
Allometry
y = kmb
(1)
m is body mass, k is the allometric constant, and b is the allometric exponent. The use of logarithms makes the equation easier to visualize—the exponent becomes the slope of a straight line when the logarithm of the variable (y) is plotted against the logarithm of body mass (m) [Eq. (2)]. Unless the allolog(y) = log(k) + b log(m)
(2)
metric exponent (b) equals 1, the ratio of y/m varies with m. The terms isometry, positive allometry, and negative allometry are sometimes used for b = 1, b > 1, and b < 1, respectively. If the variable of interest is the mass of an organ, under isometry organ mass is a fixed proportion of body mass; under positive allometry larger organisms have disproportionately large organs; and under negative allometry larger organisms have disproportionately small organs. See LOGARITHM. Proportions. The mass of the brain as a proportion of body mass from adults of different mammalian species is an example of negative allometry. Proportions vary among related species of different sizes (that is, interspecific allometry); for example, the brain is about 5% of body mass in a mouse but only 0.05% in a whale (Fig. 1). Proportions often vary within a species as well. Individuals change shape during growth, a phenomenon known as ontogenetic or growth allometry; for example, among humans the head is 25% of a baby’s length but only 15% of an adult’s height. Within species, small and large adults of the same sex may be isometric with respect to some bodily proportions, but there may also be static intraspecific allometry. The different
1 10 100 1000 10.000 body mass, kg
Fig. 1. Allometric lines for the mass of some organs across mammal species. (These are best-fit lines: individual species may lie above or below them.) Note that the skeleton is disproportionately heavy in large species (positive allometry), whereas the brain is disproportionately light (negative isometry). Heart and lungs scale approximately isometrically.
types of allometry (Fig. 2) often show a character scaling in different ways, reflecting different underlying causes. Even if small and large organisms are the same shape, size is significant because surface area and volume scale differently with body length. Consider what happens to a cube when the length of each side is doubled: the surface area increases fourfold, the volume (and hence mass) eightfold. The same relations hold for any shape: the surface area is proportional to the length squared and the mass is proportional to the length cubed. It follows, therefore, that larger organisms have less surface area per unit mass. This relationship has profound consequences, helping to explain, for instance, why shrews and hummingbirds are the smallest endotherms—any organism much smaller would lose heat too quickly through its relatively large surface. Size is also
interspecific allometry
log organ size
The study of changes in the characteristics of organisms with body size. Characteristics such as body parts or timing of reproductive events do not necessarily change in direct proportion to body size, and the ways in which they change relative to body size can often provide insights into organisms, construction and behavior. Large organisms are often not merely magnified small ones. Variables. Many characteristics, ranging from brain size and heart rate to life span and population density, change consistently with body size. These relationships normally fit a simple power function given by Eq. (1), where y is the variable under study,
static intraspecific allometry
log body size Fig. 2. Three different types of allometry. Each curve shows the growth allometry of organ size in a different species. The adults of each species show static intraspecific allometry. The circles indicate a particular stage of the life history (for example, sexual maturity). The line of interspecific allometry links points from different species.
479
480
Allosteric enzyme significant for rates; for example, longer limbs can take longer strides, so larger animals need fewer strides to cover a given distance. Across mammal species, limb length scales as m0.35, so the frequency of strides of two animals running at the same speed will differ according to their differences in m0.35. Allometric lines. In themselves, allometric lines are rough and probably simplified descriptions of the combined effects of the many (usually unknown) selective forces correlated with size; it still must be explained why the exponent takes the value that it does. For instance, in the skeletons of geometrically similar land animals, mass is proportional to length cubed while the cross-sectional area of any given bone, and hence its strength, scales as the square of its length. Beyond some critical size, these animals would collapse under their own weight. In fact, skeletal mass is proportional to m1.08 for mammals and birds because bones face static (compressing) stress when the animal is stationary and elastic (buckling) stress during locomotion. Mechanical principles can be used to ascertain whether skeletons are adapted to prevent fracture by compression or by buckling. See BONE. The leg bones of a quadruped with a vertical posture are like tall, thin, cylinders. To withstand compression equally at all body weights—a condition known as static stress similarity—their crosssectional area must increase as m1 and their length scales as m1/3, so bone mass should scale as (m1) (m1/3) = m4/3. But bone mass is proportional to m1.08, so skeletons do not show static stress similarity. Bones may instead be adapted to withstand buckling equally at all body sizes—the condition known as elastic stress similarity. The limb bones of hoofed mammals fit the model of a column, but those of other mammals do not, perhaps because their limbs are less pillarlike and so face different stresses. Elastic stress similarity is also found further afield. Limb bones of bipedal theropod dinosaurs fit the theory. Tree trunks and limbs are also thin columns faced with buckling stresses: tree trunks retain elastic stress similarity as they grow, and branch thickness scales with length as predicted. See STRESS AND STRAIN. Many allometries are less well understood. Across mammal species, for example, small mammals clearly get enough oxygen for their maximal needs, so large species would seem to have more lung than they need. Deviations from the line. Allometric lines are often used as baselines for comparison: if it is assumed that a line shows how organ mass has to change with body mass in order to keep doing the same job, organisms above or below the line have larger or smaller organs than expected for their body size. The question then arises as to why particular groups deviate from the line. Finding other variables that are correlated with this body-size-independent variation allow for powerful tests of hypotheses of adaptation. See ADAPTATION (BIOLOGY). For example, across mammal species the scaling of brain mass (proportional to m0.7) is similar to the
scaling of metabolic rate (proportional to m0.75). The similarity prompted the suggestion that brain size of newborn mammals (which is linearly related to adult brain size) scales as m0.75 because it is somehow constrained by the mother’s metabolic rate. However, another hypothesis also predicts an exponent near 0.7; surface area scales as m2/3, as discussed earlier. If processing requirements are set by the amount of peripheral sensory input, brain size might scale as m2/3 too. Choosing between these hypotheses without very good estimates of the exponent seems impossible. However, deviations from the line provide a critical test. The metabolic rate argument predicts that species with fast metabolisms for their size will also have relatively large brains; however, data do not support this prediction. Once deviations are understood, they can be used to make inferences. For instance, eye-socket size scales allometrically among extant primate species, but nocturnal species lie above the line. This knowledge leads to the inference of the activity patterns of species known only from fossils. Using allometry in this way assumes that the part under study has the same function throughout the group being considered. If function differs, so may the scaling. For example, wing length scales as m1/3 across most bird species. Hummingbird wings, however, must support the bird during hovering as well as more conventional flight, and hovering requires relatively larger wings, so hummingbird wing length scales as m2/3. Models. Allometry is perhaps most powerful when known allometric relations are coupled with theory to produce new hypotheses that can be tested. For example, mammal species with large bodies grow more slowly, mature later, have fewer but larger young after longer gestation periods, and live longer than do small species; these traits remain correlated with each other among deviations from the allometric lines. Allometric models predict not only the allometric exponents but also the correlations among the deviations. The models even provide a framework for asking questions vital to understanding scaling, such as what factors affect the evolution of body size itself. See BIOPHYSICS. Andy Purvis; Paul H. Harvey Bibliography. R. McN. Alexander, Optima for Animals, 1982; E. L. Charnov, Life History of Invariants, 1993; P. H. Harvey and J. R. Krebs, Comparing brains, Science, 249:140–146, 1990; P. H. Harvey and M. D. Pagel, The Comparative Method in Evolutionary Biology, 1991; T. A. McMahon and J. T. Bonner, On Size and Life, 1983; K. Schmidt-Nielsen, Scaling: Why Is Animal Size So Important?, 1984.
Allosteric enzyme Any one of the special bacterial enzymes involved in regulatory functions. End-product inhibition is a bacterial control mechanism whereby the end product of a biosynthetic pathway can react with the first enzyme of the pathway and prevent its activity. This end-product inhibition is a device through which
Allosteric enzyme
rate of reaction
A
B
substrate concentration Fig. 1. Rate of reaction of an enzyme in relation to substrate concentration. Curve A obtained with ordinary enzymes, curve B obtained with allosteric enzymes.
a cell conserves its economy by shutting off the synthesis of building blocks when too many are present. See BACTERIAL PHYSIOLOGY AND METABOLISM. Kinetics. Enzymes are protein substances which act as biological catalysts in converting one substance, the substrate, to another, the product. To understand the nature of the special allosteric enzymes, they must first be compared with ordinary enzymes as to their kinetics, that is, the way they work with increasing amounts of substrate. With ordinary enzymes the rate of reaction increases as the amount of substrate is increased until eventually a saturation level is obtained and no further increase occurs. The curve obtained is described as a hyperbola (Fig. 1). From this curve it can be said that the substrate molecules tend to occupy more and more binding sites until all the sites are in use. The allosteric enzymes, when similarly plotted, instead of showing a hyperbolic shape, often show a sigmoid, or S-shaped, curve. This means that there is a cooperative effect between substrate molecules so that, when the first molecules are bound to the enzyme, the subsequent ones are bound more readily. A clue to this concept of cooperative substrate interaction was obtained by the study of how the hemoglobin of the blood reacts with oxygen. Here, too, the curve describing the saturation of hemoglobin with oxygen traces a sigmoid shape. Inhibition. Another characteristic of enzymes is that their action can be prevented by inhibitors. In ordinary enzymes the inhibitors are substances which closely resemble the structure of the substrate. Because of this structural similarity, they can jam the site that the substrate normally occupies. They can thus be likened to a false coin which fits the slot of the machine but does not operate it. The surprising feature about the inhibitors of allosteric enzymes is their lack of structural similarity with the substrate. Hence, they cannot be expected to work by competing with the substrate for a common slot. Instead,
they apparently have their own special site of attachment, and when they become attached there, they somehow alter the site where the substrate fits. The substrate kinetics of the inhibited reaction supports this concept. Furthermore, S-shaped curves are also obtained when the inhibition rates are plotted against amounts of the inhibitor. This means that there is cooperative interaction between inhibitor molecules in the same way as there is cooperative interaction between substrate molecules. Activation. The initial idea of allosteric inhibition was used as an explanation for the negative feedback control shown by the regulatory enzymes. There are also times when the control must have a positive effect, and these are brought about by allosteric activators. Consider, for example, two parallel production lines of metabolic activity leading to two different end products that must be built into a common structural unit. If there is too much of one of the building blocks, then the regulatory mechanisms can react in two ways. In one case the end product can shut off its own production by negative feedback control, and in the other case it can speed up the other parallel line by positive control. In this way the activity of allosteric enzymes may be either decreased by inhibitors or increased by activators. Again, if the substrate curves are studied with and without the activator, it is found that the allosteric activator converts the S-shaped curve into the ordinary hyperbolic type. Desensitization. It is possible to treat allosteric enzymes in various ways so that they lose their sensitivity to the allosteric inhibitor. The enzymes, though desensitized with respect to inhibition, retain their activity with the substrate and, indeed, this activity is often increased to the same extent as is obtained with an activator. The S-shaped kinetics is changed to the normal variety. The agents that can cause this desensitization include heat, acid, heavy metals such as mercury, and agents which react with sulfhydryl (SH) groups on the proteins. These agents are known to be able to separate complex proteins into subunits. Allosteric enzymes are therefore considered to be complex proteins made up of identical monomers. The single unit, the monomer, is active with normal hyperbolic kinetics and is not affected by the regulatory molecules (inhibitors and activators). The multiple structure, several units combined, is the form that shows the properties of allosterism, S-shaped kinetics, and alteration by the regulators. Allosteric model. The above considerations have led to the development of a model to explain the action of allosteric enzymes (Fig. 2). It suggests that the enzyme molecule is a complex consisting of identical subunits, each of which has a site for binding the substrate and another for binding the regulatory substance. These subunits interact in such a way that two conformational forms may develop. One form is in a relaxed condition (R state) and has affinity for the substrate and activator; the other form is in a constrained or taut condition (T state) and has affinity for the inhibitor. The forms exist in a state of equilibrium, but this balance can be readily tipped
481
482
Allotheria R state
monomers
Key:
= activator
T state
= substrate
= inhibitor
= reaction site
Fig. 2. Model of allosteric enzyme composed of two subunit monomers. Regulatory effects are obtained by shifting back and forth between the two states, relaxed (R state) and taut (T state). In the R state it reacts with substrate and activators, and in the T state it reacts with inhibitors. (After J. Monod, J. Wyman, and J.-P. Changeux, On the nature of allosteric transitions: A plausible model, J. Mol. Biol., 12:88, 1965)
by binding one of the reactants. The substrate and activator are bound by the relaxed form; when this happens, the balance is tipped in favor of that state. Conversely, the inhibitor will throw the balance toward the constrained state. The balance is thus tipped one way or the other, depending on the relative concentrations of substrate and inhibitor. Since the two states require subunit interaction for their maintenance, it can be seen why dissociation of the subunits leads to a simple monomeric enzyme which no longer exhibits allosteric effects. The model also shows how the binding sites may interact in either a cooperative or antagonistic manner. See ENZYME. Joseph S. Gots Bibliography. J.-P. Changeux, Control of biochemical reactions, Sci. Amer., 212(4):36–45, 1965; B. D. Davis, Microbiology, 3d ed., 1980; J. Monod, J. Wyman, and J.-P. Changeux, On the nature of allosteric transitions: A plausible model, J. Mol. Biol., 12:88, 1965.
Allotheria One of the four subclasses of Mammalia, containing a single order, the Multituberculata. The Allotheria first appeared in the Late Jurassic and survived well into the Cenozoic, a period of at least 100,000,000 years. Fossils are known from North America, Europe, and Asia. The diagnostic features of the subclass are in the dentition, which is very specialized. There is a pair of enlarged incisors above and below; reduced lateral incisors may persist in the upper jaw. Canines are absent, leaving a diastema between incisors and cheek teeth. Premolars are variable and often reduced. The lower molars have five or more cusps in two parallel longitudinal rows, the upper molars two or three parallel rows of cusps; hence the molars are multituberculate. See DENTITION; MAMMALIA; MULTITUBERCULATA. D. Dwight Davis; Frederick S. Szalay
Allowance An intentional difference in sizes of two mating parts. With running or sliding fits, in which mating parts move relative to each other, allowance is a clearance, usually for a film of oil. In this sense, allowance is the space “allowed” for motion between parts. To avoid binding between parts, a minimum clearance is critical for a running fit. With force or shrink fits, in which mating parts, once assembled, are fixed in position relative to each other, allowance is an interference of metal; that is, a portion of metal in one part tends to occupy the same space as the adjacent portion of metal in the mating part. In this sense, allowance is the interference “allowed” to produce pressure between parts. To avoid breaking the external part, a maximum interference is critical for a drive, force, or shrink fit. Fits are classed as running or sliding fits for parts that move freely against each other. Location fits provide accurate orientation of mating parts. Force fits require appreciable assembly pressure and produce more or less permanent assembly of parts. See FORCE FIT; LOCATION FIT; PRESS FIT; RUNNING FIT; SHRINK FIT; TOLERANCE. Paul H. Black
Alloy A metal product containing two or more elements as a solid solution, as an intermetallic compound, or as a mixture of metallic phases. This article will describe alloys on the basis of their technical applications. Alloys may also be categorized and described on the basis of compositional groups. See INTERMETALLIC COMPOUNDS; METAL; SOLID SOLUTION. Except for native copper and gold, the first metals of technological importance were alloys. Bronze, an alloy of copper and tin, is appreciably harder than copper. This quality made bronze so important an alloy that it left a permanent imprint on the
Alloy civilization of several millennia ago now known as the Bronze Age. Today the tens of thousands of alloys involve almost every metallic element of the periodic table. See BRONZE; PERIODIC TABLE; PREHISTORIC TECHNOLOGY. Alloys are used because they have specific properties or production characteristics that are more attractive than those of the pure, elemental metals. For example, some alloys possess high strength, others have low melting points, others are refractory with high melting temperatures, some are especially resistant to corrosion, and others have desirable magnetic, thermal, or electrical properties. These characteristics arise from both the internal and the electronic structure of the alloy. For a discussion of the physical structures of alloys and the interatomic forces in alloys. See ALLOY STRUCTURES. Bearing alloys. These alloys are used for metals that encounter sliding contact under pressure with another surface; the steel of a rotating shaft is a common example. Most bearing alloys contain particles of a hard intermetallic compound that resist wear. These particles, however, are embedded in a matrix of softer material which adjusts to the hard particles so that the shaft is uniformly loaded over the total surface. The most familiar bearing alloy is babbitt metal, which contains 83–91% tin (Sn); the remainder is made up of equal parts of antimony (Sb) and copper (Cu), which form hard particles of the compounds SbSn and CuSn in a soft tin matrix (Fig. 1). Other bearing alloys are based on cadmium (Cd), copper, or silver (Ag). For example, an alloy of 70% copper and 30% lead (Pb) is used extensively for heavily loaded bearings. Bearings made by powder metallurgy techniques are widely used. These techniques are valuable because they permit the combination of materials which are incompatible as liquids, for example, bronze and graphite. Powder techniques
0.5 mm
Fig. 1. Babbitt metal (84% Sn–7% Cu–9% Sb). The white areas are crystals of CuSn (star dendrites) and SbSn (nearly square); these resist wear. The dark matrix is a soft, tin-rich alloy, which adjusts to the hard particles to give uniform loading. (From R. M. Brick et al., Structure and Properties of Alloys, 3d ed., McGraw-Hill, 1965)
also permit controlled porosity within the bearings so that they can be saturated with oil before being used, the so-called oilless bearings. See ANTIFRICTION BEARING; COPPER ALLOYS; POWDER METALLURGY; WEAR. Corrosion-resisting alloys. Certain alloys resist corrosion because they are noble metals. Among these alloys are the precious-metal alloys, which will be discussed separately. Other alloys resist corrosion because a protective film develops on the metal surface. This passive film is an oxide which separates the metal from the corrosive environment. Stainless steels and aluminum alloys exemplify metals with this type of protection. Stainless steels are iron alloys containing more than 12% chromium (Cr). Steels with 18% Cr and 8% nickel (Ni) are the best known and possess a high degree of resistance to many corrosive environments. Aluminum (Al) alloys gain their corrosion-deterring characteristics by the formation of a very thin surface layer of aluminum oxide (Al2O3), which is inert to many environmental liquids. This layer is intentionally thickened in commercial anodizing processes to give a more permanent Al2O3 coating. Monel, an alloy of approximately 70% nickel and 30% copper, is a well-known corrosionresisting alloy which also has high strength. Another nickel-base alloy is Inconel, which contains 14% chromium and 6% iron (Fe). The bronzes, alloys of copper and tin, also may be considered to be corrosion-resisting. See ALUMINUM ALLOYS; CORROSION; IRON ALLOYS; NICKEL ALLOYS; STAINLESS STEEL. Dental alloys. Amalgams are predominantly alloys of silver and mercury (Hg), but they may contain minor amounts of tin, copper, and zinc (Zn) for hardening purposes, for example, 33% silver, 52% mercury, 12% tin, 2% copper, and less than 1% zinc. Liquid mercury is added to a powder of a precursor alloy of the other metals. After being compacted, the mercury diffuses into the silver-base metal to give a completely solid alloy. Gold-base dental alloys are preferred over pure gold (Au) because gold is relatively soft. The most common dental gold alloy contains gold (80–90%), silver (3–12%), and copper (2–4%). For higher strengths and hardnesses, palladium and platinum (up to 3%) are added, and the copper and silver are increased so that the gold content drops to 60–70%. Vitallium [an alloy of cobalt (65%), chromium (5%), molybdenum (3%), and nickel (3%)] and other corrosion-resistant alloys are used for bridgework and special applications. See AMALGAM; DENTISTRY; GOLD ALLOYS; SILVER ALLOYS. Die-casting alloys. These alloys have melting temperatures low enough so that in the liquid form they can be injected under pressure into steel dies. Such castings are used for automotive parts and for office and household appliances which have moderately complex shapes. This processing procedure eliminates the need for expensive machining and forming operations. Most die castings are made from zincbase or aluminum-base alloys. Magnesium-base alloys also find some application when weight reduction is paramount. Low-melting alloys of lead and tin are not common because they lack the necessary
483
484
Alloy strength for the above applications. A common zincbase alloy contains approximately 4% aluminum and up to 1% copper. These additions provide a second phase in the metal to give added strength. The alloy must be free of even minor amounts (less than 100 ppm) of impurities such as lead, cadmium, or tin, because impurities increase the rate of corrosion. Common aluminum-base alloys contain 5–12% silicon (Si), which introduces hard-silicon particles into the tough aluminum matrix. Unlike zinc-base alloys, aluminum-base alloys cannot be electroplated; however, they may be burnished or coated with enamel or lacquer. Advances in high-temperature diemold materials have focused attention on the diecasting of copper-base and iron-base alloys. However, the high casting temperatures introduce costly production requirements, which must be justified on the basis of reduced machining costs. See METAL CASTING. Eutectic alloys. In certain alloy systems, a liquid of a fixed composition freezes to form a mixture of two basically different solids or phases. An alloy that undergoes this type of solidification process is called a eutectic alloy. A typical eutectic alloy is formed by combining 28.1% of copper with 71.9% of silver. A homogeneous liquid of this composition on slow cooling freezes to form a mixture of particles of nearly pure copper embedded in a matrix (background) of nearly pure silver. The advantageous mechanical properties inherent in composite materials such as plywood composed of sheets or lamellae of wood bonded together and fiber glass in which glass fibers are used to reinforce a plastic matrix have been known for many years. Attention is being given to eutectic alloys because they are basically natural composite materials. This is particularly true when they are directionally solidified so as to yield structures with parallel plates of the two phases (lamellar structure) or long fibers of one phase embedded in the other phase (fibrous structure). Directionally solidified eutectic alloys are being given serious consideration for use in fabricating jet engine turbine blades. For this purpose eutectic alloys that freeze to form tantalum carbide (TaC) fibers in a matrix of a cobalt-rich alloy have been heavily studied. See EUTECTICS; METAL MATRIX COMPOSITE. Fusible alloys. These alloys generally have melting temperatures below that of tin (450◦F or 232◦C), and in some cases as low as 120◦F (50◦C). Using eutectic compositions of metals such as lead, cadmium, bismuth, tin, antimony, and indium achieves these low melting temperatures. These alloys are used for many purposes, for example, in fusible elements in automatic sprinklers, forming and stretching dies, filler for thin-walled tubing that is being bent, and anchoring dies, punches, and parts being machined. Alloys rich in bismuth (Bi) were formerly used for type metal because these low-melting metals exhibited a slight expansion on solidification, thus replicating the font perfectly for printing and publication.
High-temperature alloys. Energy conversion is more efficient at high temperatures than at low; thus the need in power-generating plants, jet engines, and gas turbines for metals which have high strengths at high temperatures is obvious. In addition to having strength, these alloys must resist oxidation by fuel-air mixtures and by steam vapor. At temperatures up to about 1380◦F (750◦C), the austenitic stainless steels (18% Cr–8% Ni) serve well. An additional 180◦F (100◦C) may be realized if the steels also contain 3% molybdenum. Both nickel-base and cobaltbase alloys, commonly categorized as superalloys, may serve useful functions up to 2000◦F (1100◦C). Nichrome, a nickel-base alloy containing 12–15% chromium and 25% iron, is a fairly simple superalloy. More sophisticated alloys invariably contain five, six, or more components; for example, an alloy called Ren´e-41 contains approximately 19% chromium, 1.5% aluminum, 3% titanium (Ti), 11% cobalt (Co), 10% molybdenum, 3% iron, 0.1% carbon (C), 0.005% boron (B), and the balance nickel. Other alloys are equally complex. The major contributor to strength in these alloys is the solution-precipitate phase of Ni3 (Ti,Al), γ . It provides strength because it is coherent with the nickel-rich γ phase. Cobalt-base superalloys may be even more complex and generally contain carbon which combines with the tungsten and chromium to produce carbides that serve as the strengthening agent (Fig. 2). In general, the cobalt-base superalloys are more resistant to oxidation than the nickel-base alloys are, but they are not as strong. Molybdenum-base alloys have exceptionally high strength at high temperatures, but their brittleness at lower temperatures and their poor oxidation resistance at high temperatures have limited their use. However, coatings permit the use of such alloys in an oxidizing atmosphere, and they are finding increased application. A group of materials called cermets, which are mixtures of metals and compounds
50 µm
Fig. 2. Cobalt-base superalloy HS31. The dispersed phase is a carbide, (CoCrW)6C, which strengthens the metal. (Haynes Stellite Co.)
Alloy such as oxides and carbides, have high strength at high temperatures, and although their ductility is low, they have been found to be usable. One of the better-known cermets consists of a mixture of titanium carbide and nickel, the nickel acting as a binder or cement for the carbide. See CERMET. Joining alloys. Metals are bonded by three principal procedures: welding, brazing, and soldering. Welded joints melt the contact region of the adjacent metal; thus the filler material is chosen to approximate the composition of the parts being joined. Brazing and soldering alloys are chosen to provide filler metal with an appreciably lower melting point than that of the joined parts. Typically, brazing alloys melt above 750◦F (400◦C), whereas solders melt at lower temperatures. A 57% copper–42% zinc–1% tin brass is a general-purpose alloy for brazing steel and many nonferrous metals. A silicon-aluminum eutectic alloy is used for brazing aluminum, and an aluminumcontaining magnesium eutectic alloy brazes magnesium parts. The most common solders are based on lead-tin alloys. The prevalent 60% tin–40% lead solder is eutectic in composition and is used extensively for electrical circuit production, in which temperature limitations are critical. A 35% tin–65% lead alloy has a range of solidification and is thus preferred as a wiping solder by plumbers. See BRAZING; SOLDERING; WELDING AND CUTTING OF METALS. Light-metal alloys. Aluminum and magnesium (Mg), with densities of 1.6 and 1.01 oz/in.3 (2.7 and 1.75 g/cm3), respectively, are the bases for most of the light-metal alloys. Titanium (2.6 oz/in.3 or 4.5 g/cm3) may also be regarded as a light-metal alloy if comparisons are made with metals such as steel and copper. Aluminum and magnesium must be hardened to receive extensive application. Agehardening processes are used for this purpose. Typical alloys are 90% aluminum–10% magnesium, 95% aluminum–5% copper, and 90% magnesium–10% aluminum. Ternary (three-element) and more complex alloys are very important light-metal alloys because of their better properties. The aluminum-zincmagnesium system of alloys, used extensively in aircraft applications, is a prime example of one such alloy system. See ALUMINUM; MAGNESIUM; TITANIUM. Low-expansion alloys. This group of alloys includes Invar (64% iron–36% nickel), the dimensions of which do not vary over the atmospheric temperature range. It has special applications in watches and other temperature-sensitive devices. Glass-to-metal seals for electronic and related devices require a matching of the thermal-expansion characteristics of the two materials. Kovar (54% iron–29% nickel–17% cobalt) is widely used because its expansion is low enough to match that of glass. See THERMAL EXPANSION. Magnetic alloys. Soft and hard magnetic materials involve two distinct categories of alloys. The former consists of materials used for magnetic cores of transformers and motors, and must be magnetized and demagnetized easily. For alternating-current applications, silicon-ferrite is commonly used. This is an
alloy of iron containing as much as 5% silicon. The silicon has little influence on the magnetic properties of the iron, but it increases the electric resistance appreciably and thereby decreases the core loss by induced currents. A higher magnetic permeability, and therefore greater transformer efficiency, is achieved if these silicon steels are grain-oriented so that the crystal axes are closely aligned with the magnetic field. Permalloy (78.5% nickel–21.5% iron) and some comparable cobalt-base alloys have very high permeabilities at low field strengths, and thus are used in the communications industry. Ceramic ferrites, although not strictly alloys, are widely used in high-frequency applications because of their low electrical conductivity and negligible induced energy losses in the magnetic field. Permanent or hard magnets may be made from steels which are mechanically hardened, either by deformation or by quenching. Some precipitation-hardening, iron-base alloys are widely used for magnets. Typical of these are the Alnicos, for example, Alnico-4 (55% iron– 28% nickel–12% aluminum–5% cobalt). Since these alloys cannot be forged, they must be produced in the form of castings. Hard magnets are being produced from alloys of cobalt and the rare-earth type of metals. The compound RCo5, where R is samarium (Sm), lanthanum (La), cerium (Ce), and so on, has extremely high coercivity. Precious-metal alloys. In addition to their use in coins and jewelry, precious metals such as silver, gold, and the heavier platinum metals are used extensively in electrical devices in which contact resistances must remain low, in catalytic applications to aid chemical reactions, and in temperaturemeasuring devices such as resistance thermometers and thermocouples. The unit of alloy impurity is commonly expressed in karats, when each karat is 1/24 part. The most common precious-metal alloy is sterling silver (92.5% silver, with the remainder being unspecified, but usually copper). The copper is very beneficial in that it makes the alloy harder and stronger than pure silver. Yellow gold is an Au-Ag-Cu alloy with approximately a 2:1:1 ratio. White gold is an alloy which ranges from 10 to 18 karats, the remainder being additions of nickel, silver, or zinc, which change the color from yellow to white. The alloy 87% platinum (Pt)–13% rhodium (Rh), when joined with pure platinum, provides a widely used thermocouple for temperature measurements in the 1830–3000◦F (1000–1650◦C) temperature range. See GOLD; SILVER. Shape memory alloys. These alloys have a very interesting and desirable property. In a typical case, a metallic object of a given shape is cooled from a given temperature T1 to a lower temperature T2 where it is deformed so as to change its shape. Upon reheating from T2 to T1, the shape change accomplished at T2 is recovered so that the object returns to its original configuration. This thermoelastic property of the shape memory alloys is associated with the fact that they undergo a martensitic phase transformation (that is, a reversible change in crystal structure
485
486
Alloy that does not involve diffusion) when they are cooled or heated between T1 and T2. For a number of years the shape memory materials were essentially scientific curiosities. Among the first alloys shown to possess these properties was one of gold alloyed with 47.5% cadmium. Considerable attention has been given to an alloy of nickel and titanium known as nitinal. The interest in shape memory alloys has increased because it has been realized that these alloys are capable of being employed in a number of useful applications. One example is for thermostats; another is for couplings on hydraulic lines or electrical circuits. The thermoelastic properties can also be used, at least in principle, to construct heat engines that will operate over a small temperature differential and will thus be of interest in the area of energy conversion. See SHAPE MEMORY ALLOYS. Thermocouple alloys. These include Chromel, containing 90% nickel and 10% chromium, and Alumel, containing 94% nickel, 2% aluminum, 3% chromium, and 1% silicon. These two alloys together form the widely used Chromel-Alumel thermocouple, which can measure temperatures up to 2200◦F (1204◦C). Another common thermocouple alloy is constantan, consisting of 45% nickel and 55% copper. It is used to form iron-constantan and copper-constantan couples, used at lower temperatures. For precise temperature measurements and for measuring temperatures up to 3000◦F (1650◦C), thermocouples are used in which one metal is platinum and the other metal is platinum plus either 10 or 13% rhodium. See LASER ALLOYING; STEEL; THERMOCOUPLE. Lawrence H. Van Vlack; Robert E. Reed-Hill Prosthetic alloys. As discussed here, prosthetic alloys are alloys used in internal prostheses, that is, surgical implants such as artificial hips and knees. External prostheses are devices that are worn by patients outside the body; alloy selection criteria are different from those for internal prostheses. In the United States, surgeons use about 250,000 artificial hips and knees and about 30,000 dental implants per year. Alloy selection criteria for surgical implants can be stringent primarily because of biomechanical and chemical aspects of the service environment. Mechanically, an implant’s properties and shape must meet anticipated functional demands; for example, hip joint replacements are routinely subjected to cyclic forces that can be several times body weight. Therefore, intrinsic mechanical properties of an alloy, for example, elastic modulus, yield strength, fatigue strength, ultimate tensile strength, and wear resistance, must all be considered. Likewise, because the pH and ionic conditions within a living organism define a relatively hostile corrosion environment for metals, corrosion properties are an important consideration. Corrosion must be avoided not only because of alloy deterioration but also because of the possible physiological effects of harmful or even cytotoxic corrosion products that may be released into the body. (Study of the biological effects of biomaterials is a broad subject in itself, often referred to
as biocompatibility.) The corrosion resistance of all modern alloys stems primarily from strongly adherent and passivating surface oxides, such as titanium oxide (TiO2) on titanium-based alloys and chromium oxide (Cr2O3) on cobalt-based alloys. The most widely used prosthetic alloys therefore include high-strength, corrosion-resistant ferrous, cobalt-based, or titanium-based alloys: for example, cold-worked stainless steel; cast Vitallium; a wrought alloy of cobalt, nickel, chromium, molybdenum, and titanium; titanium alloyed with aluminum and vanadium; and commercial-purity titanium. Specifications for nominal alloy compositions are designated by the American Society for Testing and Materials (ASTM). Prosthetic alloys have a range of properties. Some are easier than others to fabricate into the complicated shapes dictated by anatomical constraints. Fabrication techniques include investment casting (solidifying molten metal in a mold), forging (forming metal by deformation), machining (forming by machine-shop processes, including computer-aided design and manufacturing), and hot isostatic pressing (compacting fine powders of alloy into desired shapes under heat and pressure). Cobalt-based alloys are difficult to machine and are therefore usually made by casting or hot isostatic pressing. Some newer implant designs are porous-coated, that is, they are made from the standard ASTM alloys but are coated with alloy beads or mesh applied to the surface by sintering or other methods. The rationale for such coatings is implant fixation by bone ingrowth. Some alloys are modified by nitriding or ionimplantation of surface layers of enhanced surface properties. A key point is that prosthetic alloys of identical composition can differ substantially in terms of structure and properties, depending on fabrication history. For example, the fatigue strength approximately triples for hot isostatically pressing versus as-cast Co-Cr-Mo alloy, primarily because of a much smaller grain size in the microstructure of the former. No single alloy is vastly superior to all others; existing prosthetic alloys have all been used in successful and, indeed, unsuccessful implant designs. Alloy selection is only one determinant of performance of the implanted device. See METAL, MECHANICAL PROPERTIES OF; PROSTHESIS. John Brunski Superconducting alloys. Superconductors are materials that have zero resistance to the flow of electric current at low temperatures. There are more than 6000 elements, alloys, and compounds that are known superconductors. This remarkable property of zero resistance offers unprecedented technological advances such as the generation of intense magnetic fields. Realization of these new technologies requires development of specifically designed superconducting alloys and composite conductors. An alloy of niobium and titanium (NbTi) has a great number of applications in superconductivity; it becomes superconducting at 9.5 K (critical superconducting temperature, Tc ). This alloy is preferred because of its ductility and its ability to carry large
Alloy structures amounts of current at high magnetic fields, represented by Jc (H) [where Jc is the critical current and H is a given magnetic field], and still retain its superconducting properties. Brittle compounds with intrinsically superior superconducting properties are also being developed for magnet applications. The most promising of these are compounds of niobium and tin (Nb3Sn), vanadium and gallium (V3Ga), niobium and germanium (Nb3Ge), and niobium and aluminum (Nb3Al) which have higher Tc (15 to 23 K) and higher Jc(H) than NbTi. Superconducting materials possess other unique properties such as magnetic flux quantization and magnetic-field-modulated supercurrent flow between two slightly separated superconductors. These properties form the basis for electronic applications of superconductivity such as high-speed computers or ultrasensitive magnetometers. Development of these applications began using lead or niobium (Tc of 7 K and 9 K) in bulk form, but the emphasis then was transferred to materials deposited in thin-film form. Lead-indium (PbIn) and lead-gold (PbAu) alloys are more desirable than pure lead films, as they are more stable. Improved vacuum deposition systems eventually led to the use of pure niobium films as they, in turn, were more stable than lead alloy films. Advances in thin-film synthesis techniques led to the use of the refractory compound niobium nitride (NbN) in electronic applications. This compound is very stable and possesses a higher Tc (15 K) than either lead or niobium. See REFRACTORY. Novel high-temperature superconducting materials may have revolutionary impact on superconductivity and its applications. These materials are ceramic, copper oxide–based materials that contain at least four and as many as six elements. Typical examples are yttrium-barium-copper-oxygen (Tc 93 K); bismuth-strontium-calcium-copperoxygen (Tc 110 K); and thallium-barium-calciumcopper-oxygen (Tc 125 K). These materials become superconducting at such high temperatures that refrigeration is simpler, more dependable, and less expensive. Much research and development has been done to improve the technologically important properties such as Jc(H), chemical and mechanical stability, and device-compatible processing procedures. It is anticipated that the new compounds will have a significant impact in the growing field of superconductivity. See CERAMICS; SUPERCONDUCTIVITY. D. Gubser Bibliography. G. S. Brady, H. R. Clauser, and J. A. Vaccari, Materials Handbook, 15th ed., 2002; W. D. Callister, Materials Science and Engineering: An Introduction, 6th ed., 2002; J. Frick (ed.), Woldman’s Engineering Alloys, 9th ed., 2001; J. A. Helsen and H. J. Breme, Metals as Biomaterials, 1998; D. P. Henkel and A. Pense, Structures and Properties of Engineering Materials, 5th ed., 2001; J. Mordike and P. Haasen, Physical Metallurgy, 3d ed., 1996; C. T. Lynch (ed.), Handbook of Materials Sciences, vol. 2, 1975; L. H. Van Vlack, Elements of Material Science and Engineering, 6th ed., 1989.
Alloy structures Metals in actual commercial use are almost exclusively alloys, and not pure metals, since it is possible for the designer to realize an extensive variety of physical properties in the product by varying the metallic composition of the alloy. As a case in point, commercially pure or cast iron is very brittle because of the small amount of carbon impurity always present, while the steels are much more ductile, with greater strength and better corrosion properties. In general, the highly purified single crystal of a metal is very soft and malleable, with high electrical conductivity, while the alloy is usually harder and may have a much lower conductivity. The conductivity will vary with the degree of order of the alloy, and the hardness will vary with the particular heat treatment used. For commercial applications of alloys and other information see ALLOY The basic knowledge of structural properties of alloys is still in large part empirical, and indeed, it will probably never be possible to derive formulas which will predict which metals to mix in a certain proportion and with a certain heat treatment to yield a specified property or set of properties. However, a set of rules exists which describes the qualitative behavior of certain groups of alloys. These rules are statements concerning the relative sizes of constituent atoms for alloy formation, and concerning what kinds of phases to expect in terms of the valence of the constituent atoms. The rules were discovered in a strictly empirical way, and for the most part, the present theoretical understanding of alloys consists of rudimentary theories which describe how the rules arise from the basic principles of physics. Prior to a discussion of the rules proposed by W. Hume-Rothery concerning the binary substitutional alloys, some information concerning alloy composition diagrams will be presented. Robb M. Thomson Phase Diagrams Alloys are classified as binary alloys, composed of two components; as ternary alloys, composed of three components; or as multicomponent alloys. Most commercial alloys are multicomponent. The composition of an alloy is described by giving the percentage (either by weight or by atoms) of each element in it. A phase diagram is a graphical description of the kinds and amounts of the phases that can be expected in an alloy as a function of its composition, temperature, and pressure when it has reached thermodynamic equilibrium. The phases may be liquid, vapor, or solid with various ordered and disordered crystal structures. A phase diagram does not provide information about how rapidly equilibrium can be reached; when a phase diagram is determined experimentally, it is necessary not only to find out what phases are present but also to assure that the alloy is in the stable equilibrium state. Equilibrium is reached when the Gibbs free energy of the system has reached its lowest possible value. The thermodynamic principle of minimum Gibbs energy imposes
487
488
Alloy structures 950
1700
900 liquid
1600
850 temperature, °C
1400
750
continuous solid solution
700
1300
650
1200
600
1100
550
solid solution
500 450 0
10
1000
A2B intermetallic compound
20
30
40
50
60
900 70
80
90
atomic percent, B
A
temperature, °F
1500
800
100 B
Fig. 1. Phase diagram at fixed pressure for a hypothetical alloy composed of elements A and B.
restrictions on the possible forms a phase diagram can take, and even a complicated diagram with many phases is composed of only a few geometrical features. See FREE ENERGY; PHASE EQUILIBRIUM; THERMODYNAMIC PRINCIPLES. Figure 1 shows a phase diagram for a hypothetical alloy system composed of elements A and B. As with most phase diagrams, it is drawn at fixed pressure. Hence it is a two-dimensional figure with composition and temperature as the axes. There are four equilibrium phases in the A-B system: the liquid, two solid phases with crystal structures of the pure components, and an intermetallic compound. The intermetallic compound is an orderly arrangement of two A atoms for every B atom on a crystal lattice. The other solid phases are disordered solutions of A and B atoms on different crystal lattices, one with a complete range of solubility and one with only a small solubility of B in A. The phase diagram for this system maps the melting and solidification
x = fl xl + fs xs
950
fl + fs = 1
1700 900
homogeneous liquid
850
XL 1500
XS 800
us
uid
liquid + solid
liq
dus
homogeneous solid
i
sol
750
700 0 A
10
20
30
40
50
60
70
1400
80
atomic percent, B
Fig. 2. Detail of the A-B system showing the liquidus and solidus.
90
1300 100 B
temperature, °F
1600 temperature, °C
temperatures, and the solubilities of the elements in each other and in the intermetallic compound. In the shaded regions the labeled phases are the only stable states of the alloys: in these regions the alloys are homogeneous liquid, solid solution, or A2B compound. In the other regions of the diagram, two phases of different composition can coexist in twophase, or heterogeneous, equilibrium. Two key features of the A-B diagram are the two-phase equilibrium and the three-phase reaction, and these are illustrated in more detail in Figs. 2 and 3. See SOLID SOLUTION; SOLID-STATE CHEMISTRY; THERMODYNAMIC PROCESSES. Two-phase equilibrium. Figure 2 illustrates the concept of two-phase equilibrium using the specific example of the melting behavior of the A-B system. The pure elements A and B melt completely; each has a single melting-point temperature. The alloys, however, begin melting at one temperature but do not finish melting until they reach a higher temperature. In between, the alloy is slushy, and liquid and solid coexist in equilibrium with each other. Upon cooling the liquid alloy, the first small amount of solid to precipitate as the alloy solidifies is richer in element B than the original liquid. This phenomenon is called segregation of the two constituents. The lens region of the diagram is the region of the two-phase equilibrium between liquid and solid. The upper line where the liquid begins to solidify is the liquidus, and the lower line where the solid first begins to melt is the solidus. For an alloy of overall composition x, the compositions of liquid and solid, xl and xs, in equilibrium at a particular temperature are found by reading off the compositions of the liquidus and solidus at that temperature. From the overall composition x, the amounts of liquid and solid, fl and fs, can also be calculated. For a two-phase alloy, the overall composition is obtained by averaging the compositions of liquid and solid, weighted by the amounts present. Since the overall composition is not changed during segregation, Eqs. (1) are valid. These equations are (1)
a statement of the lever rule, one of the important rules for reading phase diagrams. Liquid and solid equilibrium for fixed pressure over a range of temperatures does not occur in pure substances but only in alloys. The phase diagram can be used to predict how molten A-B alloys will solidify in different conditions. If the temperature of a liquid alloy is brought below the liquidus, a small amount of solid will form with a composition different from that of the liquid. How the solidification proceeds now depends on the rate of cooling. If the alloy is cooled very slowly through the two-phase region, equilibrium can be reached at each temperature and the compositions of the liquid and solid can continuously adjust to the values on the phase diagram. If the cooling rate is too high for equilibrium to be reached at each temperature, which is typically the
Alloy structures 1200 1100
2000 liquid
1000 temperature, °C
800
1600
liquid + (Cu) (Ag)
700 600
(Cu) 1400
eutectic reaction
liquid + (Ag)
1200 1000
(Ag) + (Cu)
500
800
400
600
300 200 0
temperature, °F
1800
900
10
20
30
40
50
60
70
80
90
400 100 Cu
weight percent copper
Ag
Fig. 3. Phase diagram of the silver-copper system, which exhibits the eutectic reaction. Silver-rich and copper-rich phases are denoted by (Ag) and (Cu), respectively.
energy is lower yet than for either a uniform solid or a uniform liquid. The only constraint on the way the alloy can split up is that the slopes of the Gibbs energies of the two phases in equilibrium must be equal. The slope of the Gibbs energy regulates the composition exchange taking place between the two phases, and at equilibrium the compositions must be constant. The compositions can be determined by drawing a common tangent to the two Gibbs energy functions; the points of tangency determine xl and xs. This method of finding a two-phase equilibrium on the phase diagram is called the common tangent construction.
Gibbs energy of the solid
Gibbs energy
case, the phase diagram can still predict what can happen. The first solid to form is rich in element B compared to the liquid. The solid does not have time to readjust its composition by diffusion of element B back out of the solid, so the liquid becomes more and more enriched in element A. The solid ends up inhomogeneous in composition, and its microstructure depends on the phase diagram and the cooling rate. The liquid-solid phase equilibria shown schematically in Fig. 2 make up the complete phase diagram for many systems, for example, silver-gold. Silver and gold are completely miscible in both the liquid and solid states and form no intermetallic compounds because of the close similarity of their chemical properties and atom sizes. Three-phase reaction. The concept of the threephase reaction is illustrated in Fig. 3 by the solidification behavior of the silver-copper alloy system. Silver and copper are not as similar to each other as are silver and gold. In the solid state, silver and copper do not mix in all proportions, and it is possible to have inhomogeneous alloys with two different solid phases in equilibrium, as well as solid phases in equilibrium with the liquid. Molten alloys can solidify to form copper-rich or silver-rich solids. The three twophases fields intersect at a horizontal line, the eutectic horizontal. If a silver-rich liquid is slowly cooled, it separates into a silver-rich solid and a copper-rich liquid. However, below the eutectic temperature, the liquid is no longer a stable phase of the system. The silver-rich solid and the copper-rich liquid react to the eutectic temperature to form a new phase, solid copper. The reaction continues at the eutectic temperature until all the liquid is consumed, leaving a two-phase mixture of silver and copper. See EUTECTICS. Minimization of Gibbs energy. The concepts of the two-phase equilibrium and the three-phase reaction describe all the ingredients of the phase diagrams of binary alloys. No more than three phases can ever coexist in equilibrium in a binary system at constant pressure. No matter how many complicated phases appear in the phase diagram, the regions are either single-phase or two-phase; horizontal lines represent the three-phase temperatures. This property of binary phase diagrams stems from the thermodynamic principle of the minimization of the Gibbs energy. In alloy systems, each phase has a Gibbs energy function of composition and temperature at fixed pressure. For a pure substance, composition does not enter the picture, but there are several Gibbs energy functions of temperature for the different phases. Whichever function is lower at the temperature of interest is the stable phase. In a binary alloy system, the composition variable provides a new way for the system to lower the Gibbs free energy. Figure 4 shows the Gibbs energy functions for liquid and solid for a hypothetical system. At the composition x, the solid phase has a lower Gibbs energy than the liquid and so is more stable. However, if the alloy segregates into a liquid of composition xl and a solid of composition xs, the Gibbs
489
Gibbs energy of the liquid
x
xS
xL
A
percent B
Fig. 4. Gibbs energies at fixed pressure and temperature as a function of concentration, and common tangent construction, for a hypothetical system.
B
490
Alloy structures The phase diagram is mapped out by determining the positions of common tangents to the Gibbs energy functions at various temperatures. If two Gibbs energy functions intersect, it is always possible to draw a common tangent. This is why the two-phase fields can extend over a range of temperature. To have three phases in equilibrium, however, three Gibbs energy functions must line up along a common tangent. This can only happen at one temperature, above and below which the Gibbs energies will move back out of line. Four phases cannot all be lined up. This is a statement of the Gibbs phase rule for binary systems. See PHASE RULE. Ternary and multicomponent systems. In a ternary system such as silver-gold-copper, the composition is represented by two variables, for example, the percentage of silver and the percentage of gold. The temperature is also still a variable. The phase diagram must be drawn in three dimensions. The Gibbs energy functions are surfaces, and the common tangent is a tangent plane rather than a line. It is possible to find a plane tangent to three surfaces over a range of temperatures, so that for ternary systems the maximum number of phases in equilibrium is four, and increases by one for each element added. The same principles of construction of the three-dimensional phase diagram work for ternary systems, with one to four phases in equilibrium in each region, but the representation of the diagrams can become complicated. Nonequilibrium processes. The state of a real alloy depends not only on the equilibrium shown on the phase diagram, but also on how far away from equilibrium the system has been brought by the temperature treatments it has been subjected to and on how close to equilibrium the system can approach in the time given it. The system always behaves in accordance with the thermodynamics and Gibbs energy functions that underlie the equilibrium phase diagram. If an equilibrium phase is not able to nucleate and grow, then a metastable equilibrium might be achieved between two other competing phases. A phase can sometimes be taken so far from its equilibrium range that it becomes very unstable and can rapidly transform to another structure by one of various nonequilibrium transformation processes. The materials scientist uses thermodynamic ideas to read from the equilibrium phase diagrams the possibilities of these processes. In short, the phase diagram is a road map that guides the materials scientist through problems encountered in alloy design, design of processing techniques, failure analysis, and behavior of alloys in performance. Moreover, phase diagrams are not used only for metallurgical systems; the behavior of ceramics and minerals, and even mixtures of polymers, is also described by phase diagrams. Joanne L. Murray; Kirit J. Bhansali Size and Valence Factors Comparison studies of a large number of alloys have led to the rules for alloy formation formulated by W. Hume-Rothery and others. These rules have recog-
nized two general types of criteria regarding alloy formation. One relates to the relative size of the atomic constituents, the other to the valence, or the electronic behavior, of the two metals. These criteria have generally been treated independently of one another, although in reality some interaction can be expected. See VALENCE. Size factor. The empirical rules for alloy formation state that, for substitutional alloys, the size of the constituents must be approximately the same, whereas one of the constituents of an interstitial alloy must be small compared to the other. Simply stated, the size factor recognizes the importance of choosing two metals which can fit together in a lattice structure. In order to define the concept of size, one must refer to the general ideas of the band theory of solids. According to this theory, the valence electrons of the metal atoms are detached from the immediate vicinity of the atom and contribute to the conduction band. In the conduction band the electrons have many of the characteristics of free particles similar to the atoms of a gas. The conduction electrons are spread over the metal, filling the interstices between the remaining ion cores of the metal atoms. The electrons of the inner shells remain tightly bound to the individual ions, however. See BAND THEORY OF SOLIDS; FREEELECTRON THEORY OF METALS. It is assumed in any discussion of atomic size that the free electrons cause a rather weak attractive force drawing the metal together, while the ion cores strongly repel one another when one core overlaps a neighbor. The strength and form of the attractive force obviously depends on the crystal type; however, the ion cores are more nearly independent of their surroundings. To the extent that the size of the atom in the metal is given by the radius of the ion core, the concept of size is well defined. However, in those cases where some of the repulsive force is contributed by the structure-dependent free electrons, the size of an atom in an alloy will vary from one alloy to another. To the extent that the size of an ion core can be rigorously defined, the alloy will contain large misfit stresses unless the ion-core sizes of the two metal atoms have a favorable ratio. In the primary or terminal solid solution parts of the composition diagram, there are two simple possibilities. The atoms of the minority category (the solute atoms) are able to supplant the solvent atoms on the solvent atom lattice, or they may fit into the open spaces of the solvent lattice. The first case is known as a substitutional alloy, while the second is called an interstitial alloy. Interstitial alloys. Examples of interstitial alloys are some of the alloys of the transition elements and, most familiarly, those of iron. The iron lattice is facecentered cubic at medium temperatures. If the facecentered lattice is considered to be made up of hard spheres of radius a0, the largest atom which can be fitted into an interstitial position (the cube center) has a radius 0.59a0. There are only four neutral atoms which have smaller radii than this value for the transition-metal group. They are hydrogen, carbon,
Alloy structures nitrogen, and boron. Carbon is actually an exception to the rule for iron, as it has a radius of 0.63a0 for the iron lattice. Actually, these four elements are not metals in their normal state, even though they do form metallic alloys with the transition metals. See CRYSTAL DEFECTS; CRYSTAL STRUCTURE; CRYSTALLOGRAPHY. Substitutional alloys. In the substitutional primary solutions the solute atom takes the place of one of the solvent atoms. In this case the size of the solvent and solute must be nearly the same. It is possible to calculate in a crude way the maximum permissible difference in size. The size difference is reflected in a lattice distortion, which is a contribution to the internal energy of the system. In general, the formation of a homogeneous alloy is possible if the change in the free energy of the alloy relative to the two separate elements is negative, as in Eq. (2). Here F is the Helmholtz free energy, U is F = U − T S
(2)
the internal energy, T is the absolute temperature, and S is the entropy. The distortion energy increases U. On the other hand, there is an entropy of mixing when the solute atoms are distributed in a random fashion on the solvent lattice. Thus, if the distortion energy is not too large to overbalance the entropy change, the alloy can form. An estimate of the distortion energy as a function of size suggests that the size of the solute atom must be within about 15% of the size of the solvent atom. This number agrees with the empirically derived Hume-Rothery rule. See ENTROPY. The substitutional alloys are the commonest type of alloy structure and have received the most study. Cu-Au is a good example of the substitutional alloy. The atomic sizes are very close, and the electronic structures of Cu and Au are very similar. As a result, this system forms a single primary solid solution system from one end of the composition diagram to the other. The only complexities are due to ordering phenomena, which will be discussed later. Valence factor. In addition to the effects of the size of the ions on the formation of alloys, the electronic structure of the atoms involved also plays a role. It is the electronic configuration of the atoms which determines why mixtures of some elements form metallic alloys, whereas some form insulating compounds. Qualitatively, it is found that alloys are formed from the atoms of the middle of the electrochemical series of the elements. The reason is that there is always a tendency for the ions of the metallic state to polarize with respect to their neighbors and form an ionic solid instead of the metallic one. If the tendency is strong, no metallic alloy state can be formed; hence, only elements from the middle of the series, where there is little change in the ionization potential from one atom to another, can form successful alloys. See ELECTROCHEMICAL SERIES. Stoichiometry. Closely related to the polarization effect is the reason why the alloys are nonstoichiometric (that is, why they form mixtures which do
not consist of small-number ratios of one element to the other) and form homogeneous phases over wide ranges of composition. For example, copper and gold mix homogeneously for any composition. This result is in striking contrast to the behavior of the ionic compounds such as sodium chloride, NaCl. See STOICHIOMETRY. In NaCl an excess Na ion is bound in the crystal with an energy less by about 1 eV than a normal ion, and this energy discrepancy is large enough so that, when the crystal becomes far from stoichiometric, the homogeneous compound is no longer formed. The excess constituent forms a separate phase, either as Na metal crystals embedded in the matrix or as bubbles of Cl2. The energy discrepancy in the case of NaCl is primarily due to the excess of charge of one sign or the other when the wrong ion is on a lattice site. Since the lattice is completely ionized, the excess charge due to the excess ion amounts to the complete ionic charge. See IONIC CRYSTALS. The situation for a metal is completely different. The charge of the ionic core of a metal is balanced by the charge of the free electrons which each atom contributes to the conduction band. Thus, in a perfect metal, the volume occupied in the lattice by each atom is neutral. If a single gold atom is placed substitutionally in a copper lattice, there will be a tendency for the charge in the gold lattice volume to become unbalanced because the copper and gold atoms have different ionization energies. Effectively, the copper and gold atoms in the metal have different affinities for the free electrons of the valence band; hence there will be a small charge imbalance induced around the gold lattice site. However, in the case of Au-Cu, the chemical valence is equal, and the amount of polarization inducted is a small fraction of a unit charge. Screening length. A further effect is the screening of the polarization around the gold atom by the conduction electrons themselves. In this case, when the region around the gold is polarized, a voltage difference is also generated. Hence, the electrons of the metal tend to rush into the affected region to even out the discrepancy. However, it is not possible to redress the balance completely, and a region of the order of the size of the gold atomic volume itself remains polarized. The size of this region is indicated by the screening length of the electrons of the metal, which are here listed for several metals in units of lattice spacing a0: Cu, 1.1; Al, 1.4; Tl, 1.0; Fe, 3.2; and Ni, 4.8. The effect of replacing a copper atom with a gold atom in Cu-Au is thus seen to be very much less than the effect of replacing a Cl ion with an Na ion in NaCl. For this reason the copper lattice can accept an unlimited number of exchanges without breaking up its structure, and the alloy forms over the entire composition range. Cu-Au is a rather special example to pick, since the ionization potentials of Cu and Au are very close, and the size factor is very favorable. The discussion, however, does illustrate the reason why these metals do not form compounds in the chemical sense.
491
492
Alloy structures Much the same reasoning would be applicable for showing why covalent crystals should also form stoichiometric compounds, where dangling covalent bonds would have much the same effect on the energy of the crystal as does an excess charge in an ionic lattice. See CHEMICAL BONDING. Electron Compounds Even though no metal alloy forms a compound in the chemical sense, there is a characteristic behavior when the alloying constituents have differing valences, and these alloys go under the name of electron compounds. In the usual alloy the intermediate region of the phase diagram is covered by several intermediate phases, with more than the one such phase illustrated in Fig. 1. The term electron compound refers to the fact that there is considerable regularity in the types of intermediate phases which are observed. Structure and electron density. It appears that in numerous cases there is a tendency for a particular type of crystal structure to correspond to a particular electron density in the conduction band. Thus, when zinc is alloyed to copper, near the 50% composition there is a phase change from the face-centered-cubic (fcc) lattice of copper to a body-centered-cubic (bcc) lattice. If it is assumed that all the valence electrons of all the atoms are contributed to the conduction band, there will be an average of 1.5 electrons per atom of the alloy in this band, since the valence of copper is 1 and the valence of zinc is 2. For the alloys Cu3Al and Cu5Sn, the crystal structure is also bcc for the same average free electron concentration. T. B. Massalsky has listed (see table) the intermediate phases which center about the characteristic electron concentrations shown. Note that there are three different compounds centered about the electron concentration 1.5 in the table. There are also two other well-developed electron compounds, one at about concentration 1.6 and one at about concentration 1.75. Also, there are two places where the hexagonal close-packed (hcp) structures occur, one at 1.5 and one at 1.75. For example, in the case of CuSi, the hcp structure is interrupted at 1.6 by a change to the phase called the γ -brass structure, which is a very complex crystal type with 52 atoms per unit cell. Brillouin zones. The striking correlations of the various electron compounds in the empirical tabulation of the table has been interpreted on the basis of the free-electron theory of metals. H. Jones has used the free-electron theory of metals modified by perturbation theory to derive an expression for the difference in energy of the electrons of the fcc and bcc crystals. His result is that as the Brillouin zone is filled, at first the energy is very nearly equal for the two crystals. The energy of the fcc crystal then drops below that of the bcc crystal at an electron concentration of about 1 electron per atom; but at a concentration close to 1.5, the situation reverses and the bcc crystal becomes more stable. Similar discussions have been given for the hexagonal lattices; these show that certain distortions of the symmetry of the per-
Binary systems with intermediate phases at concentrations of 1.5, 1.6, and 1.75 electrons per atom 1.5 Body centered cubic
Hexagonal closepacked
Manganese
Cu-Be Cu-Zn Cu-Al Cu-Ga Cu-In Cu-Si Cu-Sn Ag-Mg Ag-Zn Ag-Cd Ag-Al Ag-In Au-Mg Au-Zn Au-Cd Fe-Al Co-Al Ni-Al Ni-In Pd-In
Cu-Ag Cu-Ge Cu-Si Ag-Al Ag-In Ag-Sn Ag-Sb Ag-Cd Au-Sn Au-In
Cu-Si Ag-Hg Ag-Al Au-Al Co-Zn
1.6
1.75
γγ-
Hexagonal closepacked
Brass Cu-Zn Cu-Cd Cu-Hg Cu-Al Cu-Ga Cu-In Cu-Si Cu-Sn Ag-Zn Ag-Cd Ag-Hg Ag-In Au-Zn Au-Cd Au-In Mn-Zn Fe-Zn Co-Zn Ni-Be Ni-Zn Ni-Cd Rh-Zn Pd-Zn Pt-Be Pt-Zn Na-Pb
Cu-Zn Cu-Cd Cu-Sn Cu-Si Ag-Zn Ag-Cd Ag-Al Au-Zn Au-Cd Au-Sn Au-Al
fect hexagonal lattice which occur are due to electronic effects related to filling the Brillouin zone. See BRILLOUIN ZONE. However, in spite of the seeming success of the theory in confirming the electron compound concept in terms of the interaction of the electrons with the boundaries of the Brillouin zone, the Jones treatment must be considered a very limited theory. First of all, Jones assumed, when he adopted the freeelectron picture, that the role of the ions in the alloy is a minor one. On the other hand, when a zinc atom is placed in a copper lattice, it is necessary that the vicinity of the zinc atom be highly polarized, according to the discussion of the preceding section. In a 50% alloy of Cu-Zn, the atomic volume of a zinc ion will have only 1.5 electrons to cancel the charge of a double ionized core. Thus, the lattice cell of the zinc atom has a net charge of +1/2 . Hence, the valence electrons from zinc atoms cannot be freely contributed to the free-electron cloud of the crystal. There must be a considerable clumping of this charge around the various different ions of the alloy in such a way that the excess charge of the ions is screened. In addition, there is some evidence from the experiments on nuclear magnetic resonance in dilute alloys that the electron density in the conduction band does not change as the divalent ion is added, in direct contradiction to the simple free-electron picture. In view of these considerations, it is surprising that such simple results as those obtained by Jones are correct. In later work J. Friedel has been able to show that, although the description of the electrons in the conduction band is far from the free-electron picture
Alloy structures of Jones, the so-called rigid band model is valid for dilute alloys and leads to results similar to those of Jones. The theory of electron compounds is not yet firmly established. Order-Disorder in Alloys
Σ
In the preceding discussion on the formation of alloys, it was tacitly assumed that the atoms of the two constituents were randomly distributed on the lattice of the alloy. In very dilute primary solutions the random distribution holds because the free energy of the system in the completely random state is lower than that of a more symmetric state. To be more specific, in a dilute solution of metal B in metal A, if there is no gain of internal energy by ordering B (placing the atoms of B in a regular arrangement on the lattice sides of A), the free energy, U − TS, is a minimum when B is completely randomized. On the other hand, if an appreciable percentage of the lattice sites are occupied by B and if there is a difference in the interatomic forces between A-B nearest-neighbor atom pairs and those between A-A or B-B pairs, the situation is different. The internal energy of the alloy is then given by Eq. (3), where U = NAB EAB + NAA EAA + NBB EBB
1
(3)
NAB is the number of A-B nearest-neighbor pairs in the alloy, EAB is the energy of interaction between an A-B pair, and so on. It is assumed that the alloy atoms interact appreciably only when they are nearest neighbors. If the interaction between A-B pairs is greater than that between B-B or A-A pairs, then the state of lowest energy is that in which the number of like pairs is minimized. This state is also a completely ordered state of the crystal. For any given temperature the equilibrium state of the crystal is that for which the free energy is minimized, and a balance is struck between the contradictory tendencies of the U and TS terms of the free energy. At low temperatures the U term can be expected to predominate, and ordering will usually occur, while at higher temperatures the entropy term will predominate, and a transition to a disordered state is usual. Long- and short-range order. The ordering of a crystal is described in terms of two different order parameters. One is called local or short-range order, and the other is called long-range order. For the initial discussion of the long- and short-range order parameters, a 50% A-B alloy will be considered. It will also be supposed that the lattice structure is bcc and that the unlike pair interaction is the stronger one. The advantage of the bcc lattice for the discussion is that, in this lattice, all the neighbors of any given atom can be made unlike atoms. In the completely ordered state every atom of the crystal is surrounded by unlike nearest neighbors, and the crystal becomes disordered when some of the A atoms exchange with B atoms. The long-range order parameter is a measure of the number of such interchanges and expresses the number of sites
Tc
0 temperature
Fig. 5. Long-range order Σ of a crystal as a function of temperature. The order parameter has a rather precipitous drop at the transition temperature and is zero above that point.
occupied by the wrong atom, as in Eq. (4). Here na =
2na −1 N
(4)
is the number of A sites occupied by A atoms, and N is the total number of A lattice sites. The quantity varies from −1 to +1, and the state of complete disorder corresponds to = 0. The short-range order expresses the fact that the neighbors surrounding any given A atom will have a tendency to be all B atoms. If this tendency is averaged over the lattice, the short-range order σ is defined in Eq. (5). Here σ is the local-order paramσ = 2(q − 1 /2 )
(5)
eter, and q is the number of A-B pairs in the crystal divided by the total number of atom pairs in the crystal. The parameter σ has the range −1 to +1, with complete disorder at σ = 0. Variation with temperature. A good qualitative notion of the behavior of the order of a crystal is obtained by observing the variation of the order parameter of a crystal as a function of temperature. Theory predicts the curves shown in Figs. 5 and 6. The longrange order decreases from a state of complete order at absolute zero temperature. At a critical transition
1
σ
0
Tc temperature
Fig. 6. Short-range order σ as a function of temperature. It varies more slowly than long-range order. Considerable order still exists above the transition temperature.
493
494
Allspice temperature which depends on the strength of the A-B bonds relative to the A-A bonds, the long-range order drops to zero. The local order also starts at absolute zero with a maximum value, but decreases more slowly than the long-range order does and remains finite for all temperatures. The discussion up till now has been based on the model of a 50% alloy in a simple lattice. The theory for a material in which the composition is allowed to vary over the entire composition diagram is exceedingly complicated. The complication is due to the difficulty of specifying the types of order configurations which are possible, and computing their entropy. However, the general results are still comparable with the simple case. Long- and short-range order parameters can still be defined, with a transition temperature at the point where the long-range order disappears. The short-range order again persists to a considerable degree even above the transition temperature. Transition temperature. The existence of the transition temperature amounts to a change of phase, because there is an attendant singularity in the specificheat curve of the material. The phase change may be of the first or second order, depending on the type of singularity present in the long-range order at the transition temperature. The curve of Fig. 5 corresponds to a second-order transition; however, if there is a discontinuous drop of the order to zero at the transition point, the transition is first order. A general rule seems to be that for lattices in which the completely ordered state has only neighbors of A-B bonds, the transition is second order. The bcc lattice is such a case, as has already been explained. However, in close-packed lattices such as the fcc lattice, it is not possible for all the bonds of a 50% alloy to be of the sort A-B, and in such cases, it is found experimentally that the transition is first order. See PHASE TRANSITIONS. Detection of order. The presence of order in a crystal and the transition from the ordered to the disordered state are detected by a variety of techniques. The resistance of the alloy at low temperatures varies with the order of the crystal because the electron waves of the crystal are sensitive to irregularities in the crystal, so that as disorder increases, the resistivity increases also. X-ray and neutron scattering of the lattice are also functions of the degree of order for the same reason. The specific heat of the crystal varies with the order, since as the crystal loses its order, energy must be supplied to the lattice to form the A-A and B-B bonds which have higher energy. At the transition temperature the specific heat rises to a sharp peak, which is easily detected. See NEUTRON DIFFRACTION; SPECIFIC HEAT OF SOLIDS; X-RAY DIFFRACTION. Superlattices. The x-ray and neutron scattering of a completely ordered crystal have an interesting peculiarity. In the Bragg scattering processes new lines appear, called superlattice lines. The name is derived from the fact that the lines correspond to the appearance of a secondary crystalline structure in the lattice. Their explanation is clear in the light of the
discussion of the ordered lattice. For a 50% A-B lattice in the bcc form, the A atoms form the corners of the cube, while the B atoms are at the center positions. The crystal is then composed of two interpenetrating simple cubic crystals, and is called a superlattice. The x-ray lines for two simple cubic structures appear instead of those for a bcc crystal. See INTERMETALLIC COMPOUNDS. Robb M. Thomson Bibliography. W. Hume-Rothery and R. E. Smallman, The Structure of Metals and Alloys, 1988; A. G. Khachaturyan, Theory of Structural Transformations in Solids, 1983; W. G. Moffatt, Handbook of Binary Phase Diagrams, 5 vols., 1981; W. F. Smith, Structure and Properties of Engineering Alloys, 2d ed., 1992.
Allspice The dried, unripe fruits of a small, tropical, evergreen tree, Pimenta officinalis, of the myrtle family (Myrtaceae). This species (see illus.) is a native of the
(a)
(b) Allspice. (a) Branch with fruit. (b) Flowers.
West Indies and parts of Central and South America. The spice, alone or in mixtures, is much used in sausages, pickles, sauces, and soups. The extracted oil is used for flavoring and in perfumery. Allspice is so named because its flavor resembles that of a combination of cloves, cinnamon, and nutmeg. See MYRTALES; SPICE AND FLAVORING. Perry D. Strausbaugh; Earl L. Core
Almanac
Almanac A book that contains astronomical or meteorological data arranged according to days, weeks, and months of a given year and may also include diverse information of a nonastronomical character. This article is restricted to astronomical and navigational almanacs. Development. The earliest known almanac material was computed by the Egyptians in A.D. 467, and astronomical ephemerides appeared irregularly thereafter. The first regular almanac appears to have been produced in Germany from 1475 to 1531. The first almanac to be published regularly by a national government was the Connaissance des Temps (1679) by France. In 1767 the British began publishing the Nautical Almanac to improve the art of navigation. The Berliner Astronomisches Jahrbuch was first published in Germany for 1776, and the Efemerides Astronomicas was published in Spain in 1791. In the United States The American Ephemeris and Nautical Almanac has been published since 1855. Beginning with the issue for 1981, the two series of publications titled The Astronomical Ephemeris, which had previously replaced The Nautical Almanac and Astronomical Ephemeris, and The American Ephemeris and Nautical Almanac were continued with the title of The Astronomical Almanac. Over the years cooperation has developed between the organizations responsible for the preparation and publication of the almanacs. Under the auspices of the International Astronomical Union (IAU), agreements are reached concerning the bases and constants to be used in the publications and the exchange of data in different forms. Astronomical Almanac. The Astronomical Almanac contains ephemerides, which are tabulations, at regular time intervals, of the orbital positions and rotational orientation of the Sun, Moon, planets, satellites, and some minor planets. It also contains mean places of stars, quasars, pulsars, galaxies, and radio sources, and the times for astronomical phenomena such as eclipses, conjunctions, occultations, sunrise, sunset, twilight, moonrise, and moonset. This volume contains the fundamental astronomical data needed by astronomers, geodesists, navigators, surveyors, and space scientists. The theory and methods on which The Astronomical Almanac is based are provided in the Explanatory Supplement to the Astronomical Almanac. See ASTRONOMICAL COORDINATE SYSTEMS; EPHEMERIS. Navigational almanacs. While The Astronomical Almanac is basically designed for the determination of positions of astronomical objects as observed from the Earth, The Nautical Almanac and The Air Almanac are designed to determine the navigator’s position from the tabulated position of the celestial object. The ground point of the celestial object is used with the North and South poles and the observer’s assumed position to establish a spherical triangle known as the navigational triangle. By means of spherical trigonometry or navigational tables, the
navigator determines a computed altitude and the direction of the ground point (azimuth) of the celestial object with respect to his or her position. The computed altitude and azimuth enable the navigator to plot a line of position that goes through or near this assumed position. The computed altitude is compared to the observed altitude of the celestial object. From the combination of two or more observations the navigator is able to determine position. See CELESTIAL NAVIGATION. The Nautical Almanac contains hourly values of the Greenwich hour angle and declination of the Sun, Moon, Venus, Mars, Jupiter, and Saturn and the sidereal hour angle and declination of 57 stars for every third day. Monthly apparent positions are tabulated for an additional 173 navigational stars. The positions are tabulated to an angular accuracy of 0.1 minute of arc, which is equivalent to 0.1 nautical mile (0.2 km). Since tabular quantities must be combined to derive the navigational fix, this tabular accuracy is sufficient to produce a computed position with an error no greater than 0.3 to 0.4 nmi (0.6 to 0.7 km).The Nautical Almanac also contains the times of sunrise, sunset, moonrise, moonset, and twilight for various latitudes. A diary of astronomical phenomena, predictions of lunar and solar eclipses, visibility of planets, and other information of interest to the navigator are also included. The Air Almanac is arranged with two pages of data back to back on a single sheet for each day; the positions of the Sun, first point of Aries, three planets, and the Moon are given at 10-min intervals. As necessary, information is adjusted so that the tabulated data at any given time can be used during the interval to the next entry, without interpolation, to an accuracy sufficient for practical air navigation. The times of sunrise, sunset, twilight, moonrise, and moonset are also given daily. Also provided are star recognition charts; a sky diagram; rising, setting, and depression graphs; standard times; and correction tables necessary for air navigation. While designed for air navigators, The Air Almanac is used by mariners who accept the reduced accuracy in exchange for its greater convenience compared with The Nautical Almanac. A number of countries publish in their languages air and nautical almanacs that are based on, or similar to, the English language versions. Other almanacs. The Astronomical Phenomena, published annually, contains the times for phenomena such as eclipses, conjunctions, occultations, sunrise, sunset, moonrise, moonset, and moon phases, seasons, and visibility of planets. Also, the dates of civil and religious holidays are included. For surveyors, the Royal Greenwich Observatory publishes The Star Almanac, designed to permit the surveyor to determine a geographical position from celestial observations. Ephemerides of Minor Planets is prepared annually by the Institute of Applied Astronomy, St. Petersburg, and published by the Academy of Sciences of Russia. This volume contains the elements, opposition dates, and opposition ephemerides of all numbered minor planets. See ASTEROID.
495
496
Almond Computer-based almanacs. The Almanac for Computers was published annually by the U.S. Naval Observatory from 1977 to 1991 to provide the coefficients for power series and Chebyshev polynomials for computation of any specific kind of astronomical data. The Royal Greenwich Observatory began publishing in 1981 Compact Data for Navigation and Astronomy. The Naval Observatory introduced the Floppy Almanac as an annual source of astronomical data by means of widely used personal-computer software. The latest source of high-precision astronomical data is the Multi-Year Interactive Computer Almanac (MICA), which provides data on compact disks. Thus, in addition to being able to compute the almanac data as published, the user is able to compute data for a particular location and time. The Naval Observatory also prepares a Satellite Almanac, which provides accurate positions for the planetary satellites. The Naval Observatory provides software that can be used for computation of apparent places and transformation of coordinates. Astronomical data in the form of star catalogs are available from astronomical data centers in Strasbourg, France, and Goddard Space Flight Center in Maryland. Publication and distribution. The Astronomical Almanac, The Nautical Almanac, The Air Almanac, and Astronomical Phenomena are cooperative publications of the U.S. Naval Observatory and Rutherford Appleton Laboratory, and are available from the U.S. Government Printing Office and the Stationery Office. These publications are used in many countries, and similar references, generally based on the same basic data, are published in various languages by the Spanish, Chinese, Russian, Japanese, German, French, Indian, Argentinian, Brazilian, Danish, Greek, Indonesian, Italian, Korean, Mexican, Norwegian, Peruvian, Philippine, and Swedish governments. Examples of publications comparable to The Astronomical Almanac are the Astronomical Almanac by Russia, the Efemerides Astronomicas by Spain, the Chinese Astronomical Ephemeris, the Indian Ephemeris, and the Japanese Ephemeris. P. K. Seidelmann
tensive forests and thickets. The cultivated almond species apparently originated in central Asia (Iran) from hybridization among native species followed by local seedling selection by native peoples. Its cultivation spread with civilization along the shores of the Mediterranean Sea into North Africa and to Italy, Spain, southern France, and Portugal. From there seeds and scions were used to introduce the crop to other parts of the world, including the United States, Australia, South Africa, and South America. Production. In the United States, commercial production is limited to California; one-half or more of this production is exported. Spain is the second leading producer, but the amount produced is about onehalf that of the United States. Italy has historically been a leading producer, primarily from the Bari and Sicily areas, but production has declined sharply. Cultivation. Because of the early spring flowering habit, commercial culture of the almond is limited to areas with mild winters and few late spring frosts. The trees thrive under moderate to high summer temperatures. Trees can withstand considerable drought due to natural adaptation to arid conditions. Thus traditional older culture in European and Asian areas has been largely restricted to hillside and otherwise marginal areas without irrigation. However, the almond responds so well to irrigation, improved nutrition, and intensification of culture (two- to threefold increases in yield) that modern culture systems in California, and now in other parts of the world, utilize the best soil areas, supplemental irrigation, fertilization, disease and insect control, herbicide application, and all other aspects of advanced agriculture. Original establishment of almonds was by seeds, but modern cultivation includes selection of cultivars budded or grafted to rootstocks. Many cultivars exist, and in general these have arisen as local or chance seedlings from a particular production area, such as California, Spain, Sicily, Bari, and Australia.
Almond A small deciduous tree, Prunus amygdalus (also known as P. dulcis or Amygdalus communis), closely related to the peach and other stone fruits and grown widely for its edible seeds. The fruit, classified botanically as a drupe (see illus.), is analogous to that of the peach except that the outer fleshy (mesocarp) layer does not enlarge in size during the latter part of the fruit growth-curve. Instead it splits (dehisces) as it nears maturity and separates from the shell (endocarp) and then dries. The nut with the kernel, or seed, within is then released. See PEACH; ROSALES. Origin. Numerous wild almond species are found in the mountainous and desert regions extending from southwestern Europe to Afghanistan, Turkistan, and western China, sometimes occurring as ex-
Nut characteristics of the two major almond cultivars of California: (a) nonpareil, (b) mission. Left: Mature nut with dehisced hull; middle: nut with shell; right: kernels.
Alpha Centauri Most cultivars are self-incompatible, and orchards require provision for combination of cultivars to provide cross pollination. Pollinizer insects, primarily honeybees, are required to accomplish cross pollination. Some self-fertile cultivars have been discovered or produced by plant breeding. See BREEDING (PLANT). Harvesting. In California, essentially all almond orchards are mechanical-harvested. Nuts are knocked to the ground when they (particularly in the center of the tree) are dry. They are then swept into windrows and taken by machines after careful soil preparation. Hulls are removed from the nuts by machines in a central huller and delivered to the processor and handler. There nuts are fumigated for worm control and stored. Most nuts are shelled, except for a few special varieties. Kernels are sorted electronically and graded into various shapes and sizes. See AGRICULTURAL MACHINERY. In some parts of the world where culture is less intensive, almonds are knocked to the ground by mallets or poles onto canvases spread under the trees. Nuts are then collected by hand into boxes or bags, or may be lifted onto almond “boats” or “sleds.” Use. Almonds are used in a variety of products. Some are roasted whole and salted to be used as snacks. Others are blanched (the skin is removed) by steam and subjected to slicing, dicing, or halving. These may be roasted, and go into products such as candy bars, bakery products, ice cream, and almond paste, among many other uses. Almonds are of two general types: the bitter type is a source of prussic acid and flavoring extracts, and the sweet type has various food uses as described above. The almond kernels contain approximately 50% fat or oil, 20% protein, 20% carbohydrate, and a variety of minerals and vitamins. See NUT CROP CULTURE. D. E. Kester
Alpaca A member of the camel family, Camelidae, which has been domesticated for more than 2000 years. The Camelidae belong to the mammalian order Artiodactyla, the even-toed ungulates. The alpaca (Lama pacos), more restricted in its range than other species of this family, is found at elevations above 12,000 ft (3600 m) along the shores of Lake Titicaca on the boundaries of Peru and Bolivia (see illustration). As in other members of this family, the alpaca’s neck and head are elongate and the upper lip has a deep cleft. The long, slender legs terminate in two toes; the feet are digitigrade, that is, the animals walk on the toes and not on the entire foot or the tip of the digits. Although the majority of the artiodactyl animals characteristically possess horns, the alpaca and other members of this family do not. The alpaca has 36 teeth, dental formula I 1/3 C 1/1 Pm 3/3 M 3/3. The long, fine repellent hair, or wool, ranges from black to white and is highly prized for manufacturing
Alpacas (Lama pacos). (Courtesy of Brent Huffman/Ultimate Ungulate Images)
cloth, particularly the white wool. Like many breeds of domesticated animals, the alpaca has been bred to produce pure strains for the wool. Although the alpaca is raised chiefly for its wool, its flesh is palatable. The known extinct and extant species of the camel family provide an interesting example of discontinuous distribution. The family originated in North America and migrated centuries ago, so that the living species are now widely separated geographically. Among the living species, affinities and relationships are further complicated and obscured by their long history of domestication. See ARTIODACTYLA; BREEDING (ANIMAL); CAMEL; DENTITION; LLAMA; MAMMALIA; NATURAL FIBER. Charles B. Curtin Bibliography. R. M. Nowak, Walker’s Mammals of the World, Johns Hopkins University Press, 1999.
Alpha Centauri The third brightest star in the sky, apparent magnitude −0.3, and the brightest in the southern constellation Centaurus. It is the closest star to the Sun at a distance of 1.35 parsecs 2.59 × 1013 mi (or 4.16 × 1013 km), and its light takes more than 4 years to reach the Earth. It has an unusually large proper motion across the sky of 3.7 seconds of arc per year. See CENTAURUS. Alpha Centauri is in reality a triple system, the two main components orbiting each other with period of nearly 80 years. They were discovered as a telescopic visual pair in 1689, and have now been followed for more than three complete revolutions. The orbit is eccentric, and their mean separation is approximately 23 astronomical units. α Cen A and B, as they are also known, are main-sequence stars of spectral types G2 and K1, respectively, and the brightest of the two is very similar to the Sun in mass (A = 1.16
497
498
Alpha fetoprotein and B = 0.97 solar masses) luminosity, and effective temperature. Chemical analyses indicate that the relative abundance of elements heavier than helium in α Cen A and B is almost twice as large as in the Sun. See ASTRONOMICAL UNIT; BINARY STAR; SPECTRAL TYPE. The third component of this system, known as Proxima Centauri, is a faint reddish star of 11th magnitude, approximately 2.2◦ away in the sky. It was discovered in 1915, and has nearly the same proper motion and distance as the other two stars. Proxima is actually slightly closer to the Sun, which makes it the Earth’s nearest neighbor in space. The linear separation from α Cen A and B is estimated to be about 13,000 astronomical units. Although the probability of such a close association in space and in projected motion occurring by chance among field stars is very small, measurements of the velocity of Proxima relative to the close binary cannot yet determine whether it is gravitationally bound. If its velocity exceeds the escape velocity of the system at Proxima’s distance from the center of mass, it cannot be in orbit and is merely traveling together with the other stars in space. However, if it is indeed dynamically bound, the orbital period is probably on the order of 106 years. See STAR. David W. Latham
Alpha fetoprotein A glycoprotein that is normally present in significant amounts only in the serum of the fetus. It is produced in the yolk sac, the liver, and other tissues of the gastrointestinal tract. Its role is unknown, but alpha fetoprotein may function as a carrier (or modulator of the concentration) of a small ligand, as an immunosuppressive, as a modulator of intracellular transport of unsaturated fatty acids, as a factor in estrogen transport, or as a means of binding retinoic acid. Levels. Peak concentration of alpha fetoprotein in human fetal serum occurs at 13 weeks of gestation, and at that time it also reaches maximum levels in the amniotic fluid. Concentration of alpha fetoprotein in maternal serum peaks at approximately week 30–32 of gestation, reflecting the effect of both the production rate in the fetus, which is already declining at that age, and the mass of the growing fetal liver. In pregnant women, there are large differences in the concentration of maternal serum alpha fetoprotein among individuals. The differences are attributable to the inbred differences in the production of alpha fetoprotein among fetuses, the sex of the fetus (males produce more alpha fetoprotein), the permeability of the intervening tissues, the weight of the mother (which reflects her blood volume), and maternal diabetes mellitus (which lowers the concentration). Ethnic differences in maternal serum alpha fetoprotein concentrations are most distinguishable between whites and blacks, the latter having a 15% higher concentration. There is a difference of only a few percent, if that, between whites and Asian or Indian populations.
Abnormal levels. Substantially increased levels of alpha fetoprotein were first observed in association with certain tumors, especially of the liver. Subsequently, increased concentrations were noted in the amniotic fluid and maternal serum of pregnant women carrying a fetus affected by an open-neuraltube defect. The two major types of neural tube defects are anencephaly, which results in a failure of the forebrain to form, and spina bifida, which is a failure of a variable part of the spinal cord to close. Both defects result in abnormalities in the overlying skin that allow excessive amounts of serum to leak into the amniotic fluid, and then into maternal blood. Thus, screening of maternal serum is now routine. While the majority of pregnant women having a high concentration of serum alpha fetoprotein experience normal pregnancies, others manifest various problems, including multiple fetuses and placental abnormalities, or defects, such as gastroschisis, strictures of the fetal gastrointestinal or genitourinary tract, the Finnish type of nephrosis, occasional fetal tumors, and skin defects; these may result in fetal death. Further, about one out of three pregnancies in which alpha fetoprotein is increased has an adverse outcome, including fetal growth retardation, low birth weight, and increased perinatal mortality. In pregnancies where the fetus is affected by trisomy 21 (Down syndrome), there is lower concentration of alpha fetoprotein than usual. See DOWN SYNDROME. Other genetic defects, including tyrosinemia, are also associated with increased alpha fetoprotein concentration. More frequently, acquired illness is responsible for the elevation. The most common cause of slight-to-moderate increases in serum alpha fetoprotein concentration is cirrhosis of the liver; the resumed production of alpha fetoprotein occurs in regenerating liver tissue. Higher concentrations are also associated with tumors, especially of the liver, gonads, and pancreas. A. Baumgarten Bibliography. G. J. Mizejewski and H. I. Jacobson (eds.), Biological Activities of Alpha-Fetoproteins, vol. 1, 1989.
Alpha particles Helium nuclei, which are abundant throughout the universe both as radioactive-decay products and as key participants in stellar fusion reactions. Alpha particles can also be generated in the laboratory, either by ionizing helium or from nuclear reactions. They expend their energy rapidly as they pass through matter, primarily by taking part in ionization processes, and consequently have short penetration ranges. Numerous technological applications of alpha particles can be found in fields as diverse as medicine, space exploration, and geology. Alpha particles are also major factors in the health concerns associated with nuclear waste and other radiation hazards. The helium nucleus, or alpha particle (α), with mass 4.00150 atomic mass units (u) and charge +2,
Alpha particles is a strongly bound cluster of two protons (p) and two neutrons (n). Its stability is evident from massenergy conservation in the hypothetical fusion reaction (1). The product mass (=4.00150 u) is less 2p + 2n →
(1)
than the reactant mass (=2 × 1.00728 u + 2 × 1.00866 u) by 0.03038 u. By using Einstein’s relation E = mc2 (where c is the speed of light), this decrease in mass m (the alpha-particle binding energy) is equivalent to 28.3 MeV of energy E. The enormous magnitude of this energy is reflected in the fact that the fusion transformation of hydrogen into helium is the main process responsible for the Sun’s energy. See CONSERVATION OF ENERGY; ENERGY; HELIUM; NUCLEAR BINDING ENERGY; NUCLEAR FUSION; PROTON-PROTON CHAIN; STELLAR EVOLUTION. Discovery. Early absorption and deflection experiments with so-called rays from radioactive minerals revealed three distinct components, of which the least penetrating (named alpha rays by E. Rutherford in 1899) was shown to consist of positively charged particles, with a charge-to-mass ratio about half that of hydrogen. These particles were shown to be helium nuclei when helium gas was identified spectroscopically in a thin-walled glass vessel within which the rays had been trapped. See BETA PARTICLES; GAMMA RAYS. Noting that a few alpha particles were strongly backscattered by thin foils while most passed directly through, Rutherford concluded that the occasional deflections were caused by coulombic repulsion between the alpha particles and compact positive entities which occupied only a small fraction of the total volume. These observations led Rutherford to propose his so-called planetary model of the atom in 1911. With subsequent modifications and refinements, Rutherford’s visual concept of atomic structure has evolved to become the model of the atom accepted today. See ATOMIC STRUCTURE AND SPECTRA. Alpha radioactivity. Coulombic repulsion between the protons within a nucleus leads to increasingly larger ratios of neutron number N to proton number Z for stable nuclei, as the mass numbers A (= Z + N) increase. Neutron-deficient nuclei can improve their N/Z ratios by means of alpha decay. The decay occurs because the parent nucleus has a total mass greater than the sum of the masses of the daughter nucleus and the alpha particle. The energy converted from mass energy to kinetic energy is called the Q value and is given by the Einstein mass-energy equivalency relation, Eq. (2), where c is the velocity of light in vacQ = [M(Parent) − M(Daughter) − M(4 He)]c2
(2)
uum and the M is atomic mass. If the masses are given in atomic mass units, then the c2 can be replaced with the conversion factor 931.501 MeV/u. The kinetic energy released is shared between the daughter nucleus and the alpha particle in accordance with the
conservation of momentum, thus giving Eq. (3) for
M(4 (He) KE = Q 1 − M(Daughter)
(3)
the kinetic energy of the alpha particle, KEα, in terms of the Q value and the masses. Thus, each radioactive alpha-emitting nuclide emits the alpha with a characteristic kinetic energy, which is one fingerprint in identification of the emitter. See NUCLEAR REACTION. For example, in reaction (4), where the notation 235 92 U143
4 → 231 90 Th141 + 2 He2
(4)
A
Z UN, and so forth, is used, an alpha particle (that is, a
4
2He2
nucleus) is emitted from uranium-235, leaving thorium-231. For this decay, Q = (235.0439231 − 231.0362971 − 4.0026032)(931.501) MeV = 4.679 MeV, and thus the energy of the alpha is KEα = 4.679 MeV (1 − 4/231) = 4.598 MeV. An alpha particle interacts with the nucleus through the strong nuclear force, which has a very short range (∼10−15 m) and is attractive, and the repulsive electrical (Coulomb) force. In alpha decay the Coulomb force is primarily responsible for accelerating the alpha to its final energy, but before that can happen the alpha must escape from the hold of the nuclear force. Relative to a potential energy of 0 at large distances, the alpha particle inside a heavy nucleus has a large positive component (∼50 MeV) due to the Coulomb force, but a much larger negative component (∼−100 MeV) due to the nuclear force. Just outside the nuclear surface, the part due to the nuclear force drops to zero, while that due to the Coulomb force is still large (∼25 MeV). The total energy (potential and kinetic) of the typical decay alpha particle (4–10 MeV) is smaller than the potential energy at distances slightly larger than the nuclear radius. Thus, the alpha particle must penetrate this potential energy barrier (Coulomb barrier) in order to escape from the nucleus. Such barrier penetration is forbidden in classical physics but is permitted in quantum-mechanical theory. In quantum mechanics there is a straightforward procedure to calculate the probability that the alpha can penetrate the Coulomb barrier; the penetration probability is sensitively dependent on the energy of the alpha and the height of the barrier. Other factors that influence the halflife are the orbital angular momentum of the emitted alpha particle, which in quantum theory adds to the height of the Coulomb barrier, and the similarity of the alpha-emitting nucleus and the resulting daughter nucleus. See NONRELATIVISTIC QUANTUM THEORY; RADIOACTIVITY. There are three major natural series, or chains, through which isotopes of heavy elements decay by successions of alpha decays. Within these series, and with all reaction-produced alpha emitters as well, each isotope decays with a characteristic half-life and emits alpha particles of particular energies and intensities (see table). The presence of these radioactive nuclides in nature depends upon either a continuous production mechanism, for example the interaction of cosmic
499
500
Alpha particles In addition to the study of alpha-particle emitters that appear in nature, alpha decay has provided a useful tool to study artificial nuclei, which do not exist in nature due to their short half-lives. Alpha decay is a very important decay mode for nuclei far from stability with a ratio of protons to neutrons that is too large to be stable, especially for nuclei with atomic mass greater than 150 u. Because of the ease of detecting and interpreting decay alpha particles, their observation has aided tremendously in studying these nuclei far from stability, extending the study of nuclei to the very edge of nuclear existence. The measured energy of the alpha particle can be used in Eqs. (2) and (3) to determine the mass difference of the parent and daughter nuclei, and if one of the masses is known by previous experiments, the absolute value for the other mass can be determined. The half-life, the measured time for the alpha-particle rate to be reduced by a factor of 2, is used to calculate a decay rate, which is a product of the barrier penetration probability and nuclear structure factors. By dividing out the barrier penetration probability, a number (the reduced width) is obtained which is relatively large if the parent and daughter nuclei have the same spins and parities and similar nuclear wave functions, and is smaller for decays requiring a change in spin or nuclear wave functions. Systematic measurement of reduced widths over a range of nuclei enables scientists to trace changes in nuclear structure as a function of proton and neutron number. Nuclear structure information for more than 400 nuclides has been obtained in this way (Fig. 1). In addition, fine structure peaks
Energies and intensities of alpha-particle standard sources Half-life
Eα ∗
1.06 ⫻ 1011 years 1.41 ⫻ 1010 years
2.234 ± 0.003 4.013 ± 0.003 3.954 ± 0.008 4.197 ± 0.005 4.150 ± 0.005 4.400 ± 0.002 4.6877 ± 0.0015 4.6212 ± 0.0015 5.1566 ± 0.0004 5.1438 ± 0.0008 5.30438 ± 0.00007 5.48560 ± 0.00012 5.44290 ± 0.00013 5.49907 ± 0.00020 5.4563 ± 0.0002 6.11277 ± 0.00008 6.06942 ± 0.00012 6.63257 ± 0.00005
Source 147 232
Sm Th
238
U
4.468 ⫻ 10 years
235 230
U Th
7.04 ⫻ 108 years 7.54 ⫻ 104 years
239
Pu
2.411 ⫻ 104 years
210 241
Po Am
138.38 days 432.2 years
238
Pu
87.7 years
242
Cm
162.8 days
253
Es
20.4 days
9
Iα †
77 23 77 23 55 76.3 23.4 73.2 15.1 85.2 12.8 71.6 28.3 74.0 25.0 89.8
* Alpha-particle
energy in MeV. intensity in number of alpha particles per 100 disintegrations of source. † Alpha-particle
rays with the atmosphere, or extremely long halflives of heavy radioactive nuclides produced in past cataclysmic astrophysical events, which accounts for uranium and thorium ores in the Earth. The relative abundances of uranium-238, uranium-235, and their stable final decay products in ores of heavy elements can be used to calculate the age of the ore, and presumably the age of the Earth. See GEOCHRONOMETRY.
Hs Sg Rf No Fm Cf
160
Cm Pu
155
U Th
150
Ra Rn 140
Po Hg
element
145
135
Pb 130
Pt
125
Os 120
W Hf 115
Yb Er 110
Dy Gd Ce
105
Nd Ce
100
Ba Xe
95
Te 90
Sn 50
55
60
65
70
75
80
85
neutron number Fig. 1. Locations of known alpha emitters on the chart of nuclides are shown in color. Boxes are included for all nuclides with proton number Z > 49 which have been observed. The fraction of the time that a particular nuclide decays by alpha-particle emission is indicated by the shade of the colored square. Black boxes indicate the stable nuclides.
Alpha particles 0.007 0.006 range, cm
0.005 0.004 0.003 0.002 0.001 0
1
2
3 4 5 6 7 8 alpha-particle energy, MeV
10
9
Fig. 2. Range-energy relationship for alpha particles in silicon. 1 cm = 0.4 in.
interval is seen to increase to the large, pronounced Bragg peak near the end of the path, and then to fall off to a small tail at the very end, due to range straggling. An approximate relation for the range of alpha particles in air is given by Eq. (5). Here R is the R = 0.309 × E3/2
(5)
range in centimeters (1 cm = 0.4 in.) and E is the alpha particle’s initial energy in MeV. For an absorbing medium other than air, only the coefficient 0.309 would be different. According to another relationship, known as the Bragg-Kleeman rule, the relative stopping power per atom is approximately proportional to the square root of the mass number A. In biological systems, the ionization and excitation produced by alpha particles can damage or kill cells. By rupturing chemical bonds and forming highly reactive free radicals, alpha particles can be far more destructive than other forms of radiation which interact less strongly with matter. See CHARGED PARTICLE BEAMS; RADIATION BIOLOGY; RADIATION CHEMISTRY; RADIATION INJURY (BIOLOGY). Detectors. Alpha particles are detected by means of the ionization and excitation that they produce
stopping power, ion pairs per cm
appear in the alpha-particle spectra for many of these nuclides; each such fine structure peak gives similar information about an excited state in the daughter nucleus. Interactions with matter. By virtue of their kinetic energy, double positive charge, and large mass, alpha particles follow fairly straight paths in matter, interacting strongly with atomic electrons as they slow down and stop. These electrons may be excited to higher energy states in their host atoms, or they may be ejected, forming ion pairs in which the initial host atom becomes positively charged and the electron leaves. The more energetic ejected electrons, known as delta electrons, cause considerable secondary ionization, which accounts for 60–80% of the total ionization. A cascade of processes occurs along the alpha particle’s track, leading to tens of thousands of disruptive events per alpha particle. See DELTA ELECTRONS; RADIATION DAMAGE TO MATERIALS. The amount of energy expended by an alpha particle to form a single ion pair in passing through a medium is nearly independent of the alpha particle’s energy, but it depends strongly on the absorbing medium. While it takes about 35 eV in air and 43 eV in helium to form an ion pair, an energy of only 2.9 eV is required in germanium and 3.6 eV in silicon. The energies expended in gases are roughly correlated to their ionization potentials. For germanium, silicon, and other semiconductors, the lower ion pair energy is, effectively, the amount required to raise an electron to the conduction band. See IONIZATION POTENTIAL; SEMICONDUCTOR. The distance (or range) that an alpha particle travels before it stops depends both on the energy of the particle and on the absorbing medium. A 5.15-MeV alpha particle from plutonium-239 creates approximately 150,000 ion pairs, and travels about 1.5 in. (3.7 cm) in air before stopping. This same 5.15-MeV alpha particle has a range of about 0.0009 in. (0.0023 cm) in aluminum, 0.0014 in. (0.0035 cm) in water, and 0.0010 in. (0.0025 cm) in bone. The passage of alpha particles through silicon is a particularly important example (Fig. 2). The semiconductor industry now produces chips so small that alpha particles from contaminants in the packaging materials can disrupt the memory-array areas of the chips, a serious problem which has been researched in considerable detail. About 1985, successful barriers were designed that could reduce the problem significantly, but there is concern that alpha particles might be a limiting factor in the increasing miniaturization of chips. The daily number of alpha particles to which these chips are exposed is typically less than 15/in.2 (2.3/cm2). See INTEGRATED CIRCUITS; RADIATION HARDENING. As alpha particles expend energy and slow down, the frequency of their interactions per unit distance increases. The rate at which an alpha particle loses energy is called its stopping power or specific ionization. A Bragg curve shows the stopping power of an alpha particle as a function of the distance to the end of its path (Fig. 3). The energy expended per
80,000 70,000 60,000 50,000 40,000 30,000 20,000 10,000 20 18 16 14 12 10 8 6 4 residual range, cm
2
0
Fig. 3. Bragg curve for single alpha particle in air. The stopping power (specific ionization) in ion pairs per centimeter of air is plotted against residual range (distance from end of alpha particle’s path) in centimeters. 1 cm = 0.4 in.
501
502
Alpine vegetation in matter. The first detectors were photographic plates, which turned dark when alpha particles sensitized and reduced grains of silver halide in the emulsions. In magnetic spectrometers, calibrated photographic emulsions have been used to determine alpha-particle energies and intensities. In ionization chambers, the relationship between the number of ion pairs formed and the amount of charge collected is used in the detection process. In cloud chambers, tracks of droplets appear whenever alpha particles leave ion pairs to act as condensation centers. Finally, in present-day semiconductor detectors, the excitation of electrons from a valence band to the conduction band is detected electrically, and is then converted into precise energy and intensity values. See IONIZATION CHAMBER; JUNCTION DETECTOR; PARTICLE DETECTOR. Applications. In the promising medical field of charged-particle radiotherapy, alpha particles are useful in the treatment of inaccessible tumors and vascular disorders. It is clear from the Bragg curve (Fig. 3) that the ionizing power of alpha particles is concentrated near the ends of their paths. Thus they can deliver destructive energy to a tumor while doing little damage to nearby healthy tissue. With proper acceleration, positioning, and dosage, the energy can be delivered so precisely that alpha-particle radiotherapy is uniquely suited for treating highly localized tumors near sensitive normal tissue (for example, the spinal cord). See RADIOLOGY. The element-specific energies of backscattered (Rutherford-scattered) alpha particles are used in remote probes to analyze the mineral composition of geological formations. In particular, alpha particles scattered by light elements transfer more energy than those scattered by heavy elements. In another alphaparticle device, the energy from 238Pu alpha decay is reliably harnessed in batteries based on the Brayton cycle, and used to power scientific equipment left on the Moon. Large power systems of this type are contemplated for use in space stations. These examples are just a few of the widespread applications of alpha particles. See BRAYTON CYCLE; ION-SOLID INTERACTIONS; NUCLEAR BATTERY; SPACE POWER SYSTEMS. Hazards. Since an energetic alpha particle usually loses its energy by ionizing the atoms it encounters, it causes severe damage to living cells in its path. However, because of its short range, an alpha-particle source external to the body does not generally pose a problem. The major hazard for humans results if an alpha-particle emitter is ingested and becomes a longterm component of the body. Great care must be taken when working with alpha emitters to prevent such ingestion. Research on indoor air quality in homes and buildings has brought attention to the long-standing health hazard associated with radon. Small amounts of this alpha-emitting gas, which is generated in the decay of heavy elements in granite and building materials, can diffuse within buildings, to be ingested by occupants. This problem is exacerbated by the fact that increasing numbers of modern buildings are tightly sealed, concentrating the radon. Environmen-
tal sources of alpha-particle exposure also include radon in mines, polonium in cigarette smoke, and thorium in coal. See ENVIRONMENTAL RADIOACTIVITY; RADON. Along with the many other dangerous components, alpha-emitting isotopes in nuclear waste present serious health hazards. Efforts to achieve safe disposal of the growing accumulations of nuclear waste are receiving increased attention. See RADIOACTIVE WASTE MANAGEMENT. Carrol Bingham Bibliography. Y. A. Akovali, Review of alpha-decay data from doubly-even nuclei, Nucl. Data Sheets, 84:1–113, 1998; M. Eisenbud and T. F. Gesell, Environmental Radioactivity: From Natural, Industrial, and Military Sources, 4th ed., Academic Press, 1997; K. S. Krane, Introductory Nuclear Physics, Wiley, 1987; S. S. M. Wong, Introductory Nuclear Physics, 2d ed., Wiley, 1998; F. Yang and J. H. Hamilton, Modern Atomic and Nuclear Physics, McGraw-Hill, 1996.
Alpine vegetation Plant growth forms characteristic of upper reaches of forests on mountain slopes. In such an environment, trees undergo gradual changes that, though subtle at first, may become dramatic beyond the dense forest as the zone of transition leads into the nonforested zone of the alpine tundra. In varying degrees, depending on the particular mountain setting, the forest is transformed from a closed-canopy forest to one of deformed and dwarfed trees interspersed with alpine tundra species (Fig. 1). This zone of transition is referred to as the forest-alpine tundra ecotone. The trees within the ecotone are stunted, often shrublike, and do not have the symmetrical shape of most trees within the forest interior. The classic image is one of twisted, stunted, and struggling individual trees clinging to a windswept ridge (Fig. 2). The ecotone in which these trees exist is visually one of the most striking vegetational transition areas known.
Fig. 1. Forest-alpine tundra near Hyndman Cirque, Idaho. The trees change in density and form within the ecotone beyond the boundary of the timberline upslope toward the alpine tundra. The distribution of timberline, treeline, and the ecotone is irregular because of climatic, geomorphic, and topographic controls. (Courtesy of T. Crawford)
Alpine vegetation
Fig. 2. Whitebark pine (Pinus albicaulis) showing the classic dwarfed and contorted image of trees within the windswept forest-alpine tundra ecotone of the Beartooth Plateau, Wyoming. (Courtesy of T. Crawford)
These trees are often referred to as krummholz, a German term meaning crooked wood. In the correct scientific usage, krummholz refers only to those ecotonal trees of the European Alps that have inherited their deformed shape genetically. Trees that have not been proven to have the genetic deformation should be referred to as environmental dwarfs, cripples, elfin wood, or (in reference to their precise form) flag, flag-mat, and mat tree species. Whether or not the deformation is inherited, the ultimate cause is the presence of severe environmental conditions. The forest-alpine tundra ecotone is a mosaic of both tree and alpine tundra species; and it extends from timberline (the upper limit of the closedcanopy forest of symmetrically shaped, usually evergreen trees) to treeline (the uppermost limit of tree species) and the exposed alpine tundra. With elevational increases, tree deformation is magnified, tree height is reduced, and the total area occupied by trees becomes smaller as the alpine shrub, grass, and herbaceous perennials become more dominant. See PLANTS, LIFE FORMS OF. Distribution. The forest-alpine tundra ecotone is generally at its highest elevation in the tropics and lowest in the polar regions. It is also higher in continental areas than in marine locations. The precise dis-
tribution is controlled by a combination of climatic, topographic, and geomorphic processes. The effect of these active processes can be seen by an irregular distribution of timberlines and treelines (Fig. 1). Coniferous trees on many steep slopes are eliminated by avalanches; and they often are replaced by flexible deciduous species or herbaceous meadow cover, leaving trees on the ridges where, undisturbed, they may attain great ages. See ALTITUDINAL VEGETATION ZONES; PLANT GEOGRAPHY. Environmental controls. The environment in which these tenacious individuals survive is harsh and involves a complex interaction of many factors, with the major controlling factor often being climate. The climate is characterized by a short growing season, low air temperatures, frozen soils, drought, high levels of ultraviolet radiation, irregular accumulation of snow, and strong winds. The interaction of all these factors produces varying levels of stress within the trees. The ultimate cause of the tree deformations and of the eventual complete cessation of tree growth lies in the inability of the tissues of the shoots and the needles to mature and prepare for the harsh environmental conditions. As the length of the growing season decreases with elevation, new needles often do not mature; they have thinner cuticles (the waxlike covering on the needles that protects against desiccation and wind abrasion), and they are less acclimated against low air temperatures. Factors that particularly affect the length of the growing season include air and soil temperatures, and the depth and distribution of snow. Low air temperatures restrict growth by limiting net uptake and assimilation of carbon dioxide, affecting the ratio of photosynthesis to respiration. With a certain duration of heat, enzymes accumulate and initiate cell division and growth. With an adequately long growing season, the tissues of the trees ripen enough to tolerate the seasonally adverse conditions. If the season is inadequate, the tissues are not well prepared. Snow serves as the primary source of moisture for the trees; and it insulates the soil, keeping temperatures somewhat moderate. If an adequate amount of snow accumulates around the tree prior to the arrival of the low air temperatures of winter, the soil temperatures may remain above 32◦F (0◦C), providing liquid water. Availability of water for root uptake and an adequately developed cuticle are critical for preventing winter desiccation, considered a prime cause of tissue loss, deformation, and eventual death. Late-lying snow restricts the warming of soil in early summer, slowing the initiation of growth. Late-lying snow also provides a habitat for a parasitic snow fungus that kills needles covered by snow. Wind is a major factor in sculpturing the trees in the ecotone. Mechanical abrasion by windtransported snow, rock, and ice particles pits needle surfaces and breaks stems. The exposure of pitted tissue to a dry, winter atmosphere guarantees desiccation, which is usually lethal. Wind can blow snow away from the trees, exposing them to desiccation,
503
504
Alternating current low air temperatures, and high amounts of radiation. See WIND. Winter air temperatures become critical, depending on the timing and the status of the tree tissues. When adequately prepared, the tree tissues can tolerate temperatures as low as −40◦F (−40◦C). This survival mechanism, known as cold-hardy acclimation, is the seasonal transition from the tender, frost-susceptible condition to the hardy, non-frostsusceptible one. If the tissues have not adequately acclimated because of a poor growing season, the cells are lethally sensitive to temperatures above the minimums tolerated by well-prepared tissues. See ADAPTATION (BIOLOGY); AIR TEMPERATURE. Radiation is critical, because the greater amounts of ultraviolet at higher elevations are magnified by the reflective snow cover. Deactivation of the chlorophyll may occur, resulting in the weakening or death of needle tissue. Large amounts of radiation can be detrimental to photosynthesis, again restricting the growth of viable tissue tolerant of the seasonal climates. Tree forms. Changes in tree form reflect increases in severity of climate with elevational rise and topographic exposure. A common progression of form changes, upslope from timberline to treeline, is from a flag, to a flag-mat, to a mat form. The trees with a flag form have branching mainly on the leeward, protected side of the upper part of the trunk. The flag branching is exposed above the winter snow cover, and growth on the windward side is eroded away by mechanical abrasion. Snow often accumulates to great depths around the base of these trees, inducing the parasitic snow fungus there. The combination results in a distorted flag form with barren branches on the lower trunk (Fig. 3). The trees of flag-mat form support scrawny trunks that are flagged. Close to the ground, a mat, or shrublike growth, survives protected under the winter snowpack. At the highest elevations within the ecotone, the mat (cushiontree) form exists. The tree is dwarfed to a mat. Rarely are vertical branches or upright leaders found. Each mat has a size and shape that is adapted to its specific topographic environment. The wind side of the mat is a contorted mass of broken needles and distorted
branches. Survival of the mat is dependent on winter snow cover, as shoots that project above the snow are abraded mechanically. Importance of trees. The environmentally dwarfed and contorted trees of the forest-alpine tundra ecotone are important for many reasons. Snow is maintained in the ecotone by low air temperatures and tree shading, providing a water source late into the summer season. The trees afford shade, wind protection, and the moisture of snowbanks, critical for the establishment and survival of seedlings. Many of the oldest trees in the world have been found in the ecotone. These offer yearly historical records of events (such as climate changes, volcanic eruptions, pollution, and fire), as recorded in tree rings. Finally, the ecotone offers a unique habitat for other plants and animals to thrive as the two communities combine resources. Ecosystems of the forest-alpine tundra ecotone are increasingly subjected to disturbances from human activities. Large numbers of visitors enter mountain areas and forest-tundra ecotones in search of recreation and wilderness experience. The impact can be large. Vegetation is trampled, water often becomes polluted, and soil is eroded along trails, contributing to a progressive degradation of ecotones. Trees that grow only a few millimeters each year are sometimes stripped of their wood for campfires or even for decorative items. Then much time is needed to replenish the growth; and in many cases regrowth may not be possible under the present climate. Strategies for management must respond to the human use within these ecosystems, as well as to the impacts of acid rain and additional climatic variability (due to increasing amounts of carbon dioxide and other chemicals in the atmosphere and to depletion of ozone). A combination of good management and better understanding of the processes and dynamics of the forest-alpine tundra ecotone ecosystem would provide a safeguard against potential environmental degradation. See ACID RAIN; FOREST ECOSYSTEM; TREE; TUNDRA; VEGETATION AND ECOSYSTEM MAPPING. Katherine J. Hansen Bibliography. J. D. Ives (ed.), Mountains, 1994; J. D. Ives and R. G. Barry (eds.), Arctic and Alpine Environments, 1974; L. W. Price, Mountains and Man: A Study of Process and Environment, 1981; W. Tranquillini, Physiological Ecology of the Alpine Timberline, 1979.
Alternating current
Fig. 3. Flagged tree, mainly supporting branching on the leeside of the upper trunk; needles have been lost from the lower trunk to the parasitic snow fungus. (Courtesy of T. Crawford)
Electric current that reverses direction periodically, usually many times per second. Electrical energy is ordinarily generated by a public or a private utility organization and provided to a customer, whether industrial or domestic, as alternating current. See ELECTRIC POWER GENERATION. One complete period, with current flow first in one direction and then in the other, is called a cycle, and 60 cycles per second (60 Hz) is the
Alternating current customary frequency of alternation in the United States and in the rest of North America. In Europe and in many other parts of the world, 50 Hz is the standard frequency of alternation. On aircraft a higher frequency, often 400 Hz, is used to make possible lighter electrical machines. When the term alternating current is used as an adjective, it is commonly abbreviated to ac, as in ac motor. Similarly, direct current as an adjective is abbreviated dc. Advantages. The voltage of an alternating current can be changed by a transformer. This simple, inexpensive, static device permits generation of electric power at moderate voltage, efficient transmission for many miles at high voltage, and distribution and consumption at a conveniently low voltage. With direct (constant) current it is not possible to use a transformer to change voltage. On a few power lines, electrical energy is transmitted for great distances as direct current, but the electrical energy is generated as alternating current, transformed to a high voltage, rectified to direct current and transmitted, and then changed back to alternating current by an inverter, to be transformed down to a lower voltage for distribution and use. See DIRECT-CURRENT TRANSMISSION. In addition to permitting efficient transmission of energy, alternating current provides advantages in the design of generators and motors, and for some purposes gives better operating characteristics. Certain devices involving chokes and transformers could not be operated on direct current. Also, the operation of large switches (called circuit breakers) is facilitated because the instantaneous value of alternating current automatically becomes zero twice in each cycle, and an opening circuit breaker need not interrupt the current but only prevent current from starting again after its instant of zero value. See ALTERNATING-CURRENT GENERATOR; ALTERNATINGCURRENT MOTOR; CIRCUIT BREAKER. Sinusoidal form. An alternating current waveform is shown diagrammatically in Fig. 1. Time is measured horizontally (beginning at any arbitrary moment) and the current at each instant is measured vertically. In this diagram it is assumed that the current is alternating sinusoidally; that is, the current i is described by Eq. (1), where Im is the maximum i = Im sin 2πf t
(1)
instantaneous current, f is the frequency in cycles
current or voltage
1 cycle
time Fig. 1. Diagram of sinusoidal alternating current.
per second (hertz), and t is the time in seconds. See SINE WAVE. Measurement. Quantities commonly measured by ac meters and instruments are energy, power, voltage, and current. Other quantities less commonly measured are reactive volt-amperes, power factor, frequency, and demand (of energy during a given interval such as 15 min). Energy is measured on a watt-hour meter. There is usually such a meter where an electric line enters a customer’s premises. The meter may be single-phase (usual in residences) or three-phase (customary in industrial installations), and it displays on a register of dials the energy that has passed, to date, to the system beyond the meter. The customer frequently pays for energy consumed per the reading of such a meter. See ELECTRICAL ENERGY MEASUREMENT; WATTHOUR METER. Power is measured on a wattmeter. Since power is the rate of consumption of energy, the reading of the wattmeter is proportional to the rate of increase of the reading of a watt-hour meter. The same relation is expressed by saying that the reading of the watt-hour meter, which measures energy, is the integral (through time) of the reading of the wattmeter, which measures power. A wattmeter usually measures power in a single-phase circuit, although threephase wattmeters are sometimes used. See ELECTRIC POWER MEASUREMENT; WATTMETER. Current is measured by an ammeter. Power absorbed by an element is the product of current through the element, voltage across it, and the power factor between the current and the voltage, as shown in Eq. (5). With unidirectional (direct) current, the amount of current is the rate of flow of electricity; it is proportional to the number of electrons passing a specified cross section of a wire per second. This is likewise the definition of current at each instant of an alternating-current cycle, as current varies from a maximum in one direction to zero and then to a maximum in the other direction (Fig. 1). An oscilloscope will indicate instantaneous current, but the value of instantaneous current is not often useful for analysis purposes. A dc (d’Arsonval-type) ammeter will measure average current, but this is useless in an ac circuit, for the average of sinusoidal current is zero. A useful measure of alternating current is found in the ability of the current to do work, and the amount of current is correspondingly defined as the square root of the average of the square of instantaneous current, the average being taken over an integer number of cycles. This value is known as the root-mean-square (rms) or effective current. It is measured in amperes. It is a useful measure for current of any frequency. The rms value of direct current is identical to its dc value. The rms value of sinusoidally alternating current is Im/2, where Im is the maximum instantaneous current. [See Fig. 1 and Eq. (1).] See AMMETER; CURRENT MEASUREMENT; OSCILLOSCOPE. Voltage is measured by a voltmeter. Voltage is the electrical pressure. It is measured between one point and another in an electric circuit, often between the
505
506
Alternating current v = V m sin 2π ft
Whether they are sinusoidal or not, average power, rms voltage, and rms current can be measured, and a value for power factor is implicit in Eq. (5). This
i = l m sin (2π ft − ϕ) π
0
ϕ
P = V I (power factor)
2π
electrical angle = 2π ft radians
Fig. 2. Phase angle ϕ.
two wires of the circuit. As with current, instantaneous voltage in an ac circuit reverses each half cycle and the average of sinusoidal voltage is zero. Therefore the rms or effective value of voltage is used in ac systems. The rms value of sinusoidally alternating voltage is Vm/2, where Vm is the maximum instantaneous voltage. This rms voltage, together with rms current and the circuit power factor, is used to compute electric power, as in Eqs. (4) and (5). See VOLTAGE MEASUREMENT; VOLTMETER. The ordinary voltmeter is connected by wires to the two points between which voltage is to be measured, and voltage is proportional to the current that results through a very high electrical resistance within the voltmeter itself. The voltmeter, actuated by this current, is calibrated in volts. Phase difference. Phase difference is a measure of the fraction of a cycle by which one sinusoidally alternating quantity leads or lags another. Figure 2 shows a voltage waveform v which is described in Eq. (2) and a current i which is described in Eq. (3). v = Vm sin 2πf t
(2)
i = Im sin (2πf t − ϕ)
(3)
The angle ϕ is called the phase difference between the voltage and the current; this current is said to lag (behind this voltage) by the angle ϕ. It would be equally correct to say that the voltage leads the current by the phase angle ϕ. Phase difference can be expressed as a fraction of a cycle, as an angle in degrees, or as an angle in radians, as shown in Eq. (3). If there is no phase difference, and ϕ = 0, voltage and current are in phase. If the phase difference is a quarter cycle, and ϕ = ±90◦, the quantities are said to be in quadrature. Power factor. Power factor is defined in terms of the phase angle. If the rms value of sinusoidal current from a power source to a load is I and the rms value of sinusoidal voltage between the two wires connecting the power source to the load is V, the average power P passing from the source to the load is shown as Eq. (4). The cosine of the phase angle, P = V I cos ϕ
(5)
gives a definition of power factor when V and I are not sinusoidal. If the voltage across the element and the current through it are in phase (and of the same waveform), power factor equals 1. If voltage and current are out of phase, power factor is less than 1. If voltage and current are sinusoidal and in quadrature, then the power factor is zero. The phase angle and power factor of voltage and current in a circuit that supplies a load are determined by the load. Thus a load of pure resistance, such as an electric heater, has unity power factor. An inductive load, such as an induction motor, has a power factor less than 1 and the current lags behind the applied voltage. A capacitive load, such as a bank of capacitors, also has a power factor less than 1, but the current leads the voltage, and the phase angle ϕ is negative. If a load that draws lagging current (such as an induction motor) and a load that draws leading current (such as a bank of capacitors) are both connected to a source of electric power, the power factor of the two loads together can be higher than that of either one alone, and the current to the combined loads may have a smaller phase angle from the applied voltage than would currents to either of the two loads individually. Although power to the combined loads is equal to the arithmetic sum of power to the two individual loads, the total current will be less than the arithmetic sum of the two individual currents (and may, indeed, actually be less than either of the two individual currents alone). It is often practical to reduce the total incoming current by installing a bank of capacitors near an inductive load, and thus to reduce power lost in the incoming distribution lines and transformers, thereby improving efficiency. This process is called power-factor correction. Three-phase system. Three-phase systems are commonly used for generation, transmission, and distribution of electric power. A customer may be supplied with three-phase power, particularly if a large amount of power is needed or if the customer wishes to use three-phase loads. Small domestic customers are usually supplied with single-phase power. A three-phase system is essentially the same as three ordinary single-phase systems with the three voltages of the three single-phase systems out of phase with each other by one-third of a cycle (120◦), as shown in Fig. 3. The three voltages may be written van
vbn
vcn
(4)
cos ϕ, is called the power factor. Thus the rms voltage, the rms current, and the power factor are the components of power. The foregoing definition of power factor has meaning only if voltage and current are sinusoidal.
0
1 6
1 3
1 2
2 3
5 6
1
Fig. 3. Voltages of a balanced three-phase system.
cycles
Alternating current generator
load
a
an
bn
b
an
bn
cn
cn c n
Fig. 4. Connections of a simple three-phase system.
as Eqs. (6), (7), and (8), where Van(max) is the maxivan = Van(max) sin 2πf t
(6)
vbn = Vbn(max) sin 2π(f t − 1 /3 )
(7)
vcn = Vcn(max) sin 2π(f t − 2 /3 )
(8)
mum value of voltage in phase an, and so on. The three-phase system is balanced if relation (9) holds, Van(max) = Vbn(max) = Vcn(max)
(9)
and if the three phase angles are equal, one-third cycle each as shown. If a three-phase system were actually three separate single-phase systems, there would be two wires between the generator and the load of each system, requiring a total of six wires. In fact, however, a single wire can be common to all three systems, so that it is only necessary to have three wires for a threephase system (Fig. 4a–c) plus a fourth wire n serve as a common return or neutral conductor. On some systems the Earth is used as the common or neutral conductor. Each phase of a three-phase system carries current and conveys power and energy. If the three loads on the three phases of the three-phase system are identical and the voltages are balanced, then the currents are balanced also. Figure 2 can then apply to any one of the three phases. It will be recognized that the three currents in a balanced system are equal in rms (or maximum) value and that they are separated one from the other by phase angles of one-third cycle and two-thirds cycle. Thus the currents (in a balanced system) are themselves symmetrical. Note, however, that the three currents will not necessarily be in phase with their respective voltages; the corresponding voltages and currents will be in phase with each other only if the load is pure resistance and the phase angle between voltage and current is zero. Otherwise some such relation as that of Fig. 2 will apply to each phase. It is significant that, if the three currents of a threephase system are balanced, their sum is zero at every instant. In practice, the three currents are not usually exactly balanced, and one of two situations arises. Either the common neutral wire n is used, in which case it carries little current (and may be of high resistance compared to the other three line wires), or else the common neutral wire n is not used, only three line wires being installed, and the three phase currents are thereby forced to add to zero even though
this requirement results in some imbalance of phase voltages at the load. It is also significant that the total instantaneous power from generator to load is constant (does not vary with time) in a balanced, sinusoidal, three-phase system. Power in a single-phase system that has current in phase with voltage is maximum when voltage and current are maximum, and is instantaneously zero when voltage and current are zero; if the current of the single-phase system is not in phase with the voltage, the power will reverse its direction of flow during part of each half cycle. However, in a balanced three-phase system, regardless of the phase angle, the flow of power is unvarying from instant to instant. This results in smoother operation and less vibration of motors and other ac devices. Three-phase systems are almost universally used for large amounts of power. In addition to providing smooth flow of power, three-phase motors and generators are more economical than single-phase machines. Polyphase systems with two, four, or other numbers of phases are possible, but they are little used except when a large number of phases, such as 12, is desired for economical operation of a rectifier. H. H. Skilling Symmetrical (0, 1, 2) components. When three coils are equally spaced around the periphery of the stator of a generator, a properly shaped magnetic field rotating in a forward direction at uniform velocity will induce voltages in these coils. The three voltages may be written as Eqs. (6), (7), and (8). For analytical purposes it is more convenient to express the voltages as phasors, as in Eqs. (10). For a discusVa1 = Va1 Vb1 = Va1 e−j2π/3
(10)
Vc1 = Va1 e+j2π/3 sion of phasor notation and complex representation. See ALTERNATING-CURRENT CIRCUIT THEORY. In order to simplify notation, it has become accepted practice to introduce a standard set of unit three-phase phasors as in Eqs. (11), so that Eqs. (10) 1 = e j0 a = e j2π/3 a2 = e−j2π/3
(11)
may be written more succinctly in matrix format as in Eq. (12). 1 Va1 Vb1 = (Va1 ) a2 (12) Vc1 a See MATRIX THEORY. When the voltages proceed in time-phase in the order Va, Vb, Vc, as reflected in Eqs. (6), (7), and (8), the set of three-phase voltage is called a positivesequence set of voltages and is identified by the subscript 1, as in Eqs. (10) and (12). When three coils are equally spaced around the periphery of the stator of a generator, a properly shaped magnetic field rotating in a backward direction at uniform velocity will induce voltages in these coils which proceed in time-phase in the order Va, Vc, Vb, and the set of
507
508
Alternating-current circuit theory three-phase voltages is called a negative-sequence set of voltages and identified by the subscript 2. In terms of phasor notation, the voltages are written as in Eqs. (13). Va2 1 Vb2 = (Va2 ) a (13) Vc2 a2 If coils b and c are placed in the same slot with coil a, then Va = Vb = Vc, and the set of three-phase voltages is called a zero-sequence set of voltages and identified by the subscript 0. In terms of phasor notation the voltages are written as in Eqs. (14). Va0 1 Vb0 = (Va0 ) 1 (14) Vc0 1 The set of positive-, negative-, and zero-sequence components are collectively referred to as symmetrical components. If the voltages Va, Vb, and Vc are written as in Eqs. (15), and, in matrix format, as in Eq. (16), then Va = Va0 + Va1 + Va2 = Va0 + Va1 + Va2 Vb = Vb0 + Vb1 + Vb2 = Va0 + a2 Va1 + aVa2 (15) Vc = Vc0 + Vc1 + Vc2 = Va0 + aVa1 + a2 Va2 Va 1 1 1 Va0 2 Vb = 1 a a Va1 (16) 1 a a2 Vc Va2 by selecting complex values for Va0, Va1, and Va2, respectively, in an arbitrary manner, it is possible to establish an infinite number of unbalanced phase voltages Va, Vb, and Vc. Conversely, and more important, since the inverse relationship given by Eq. (17) Va 1 1 1 Va0 2 Va1 = (1/3) 1 a a Vb (17) Va2 1 a2 a Vc exists, for any given unbalanced set of three-phase voltages, it is always possible to determine a corresponding unique set of zero-, positive-, and negativesequence voltages. Equally well, for any given unbalanced set of three-phase currents, it is always possible to determine a corresponding unique set of zero-, positive-, and negative-sequence currents. Three-phase power system components, such as transmission lines, transformers, loads, motors, and generators, may, in many practical instances, be considered geometrically symmetrical among the three phases. In this event, it can be shown that when a zero-, positive-, or negative-sequence voltage is applied to such a component, the resultant current will be a zero-, positive-, or negative-sequence component, respectively. The associated zero-, positive-, and negative-sequence impedances can be readily identified. Consequently, because of the symmetrical nature of each of the symmetrical-component voltages and currents, it is necessary to consider only one of the three phases of the three-phase power system in an
analysis. By accepted convention, phase “a” is the selected phase. Perhaps the most important virtue of symmetrical components lies in the fact that if a common type of unsymmetrical fault (single-line-to-ground fault, lineto-line fault, one open conductor, two open conductors, and so forth) occurs at one point in an otherwise symmetrical three-phase network, a relatively simple interconnection occurs between the symmetrical component networks at the fault location. M. Harry Hesse Power and information. Although this article has emphasized electric power, ac circuits are also used to convey information. An information circuit, such as telephone, radio, or control, employs varying voltage, current, waveform, frequency, and phase. Efficiency is often low, the chief requirement being to convey accurate information even though little of the transmitted power reaches the receiving end. For further consideration of the transmission of information. See ELECTRICAL COMMUNICATIONS; RADIO; TELEPHONE; WAVEFORM. An ideal power circuit should provide the customer with electrical energy always available at unchanging voltage of constant waveform and frequency, the amount of current being determined by the customer’s load. High efficiency is greatly desired. See CAPACITANCE; CIRCUIT (ELECTRICITY); ELECTRIC CURRENT; ELECTRIC FILTER; ELECTRICAL IMPEDANCE; ELECTRICAL RESISTANCE; INDUCTANCE; JOULE’S LAW; NETWORK THEORY; OHM’S LAW; RESONANCE (ALTERNATING-CURRENT CIRCUITS). H. H. Skilling Bibliography. J. Arrillaga and C. P. Arnold, Computer Analysis of Power Systems, 1990; O. I. Elgerd, Electric Energy Systems Theory, 2d ed., 1982; J. J. Grainger and W. D. Stevenson, Jr., Power System Analysis, 1994; C. P. Paul, S. A. Nasar, and L. E. Unnewehr, Introduction to Electrical Engineering, 2d ed., 1992.
Alternating-current circuit theory The mathematical theory for the analysis of electric circuits when the currents and voltages are alternating functions of time. In an electric circuit carrying direct current only (Fig. 1), the current I (amperes) flowing through a circuit of total resistance R (ohms) when a steady voltage V (volts) is applied is given by Ohm’s law, Eq. (1). I=
V R
(1)
See CIRCUIT (ELECTRICITY); DIRECT-CURRENT CIRCUIT THEORY; OHM’S LAW. In many practical situations, electrical energy is generated, transmitted, and distributed in a form in which the currents and voltages are not constant but are sinusoidally varying or “alternating” functions of time (Fig. 2). Public electricity supply systems are
Alternating-current circuit theory I
of waveforms, illustrated by phasor diagrams. It will be shown how the use of complex algebra provides a more powerful and elegant theory. Waveforms. A sinusoidal (or alternating) function a of time t is shown in Fig. 2a. It is described by Eq. (2), where A0 is the maximum value (or ampli-
+
V
R
a(t) = A0 sin ωt
−
Fig. 1. Direct-current circuit.
of this type, and many telecommunications systems use electrical signals of sinusoidal form as carriers of information. It is necessary to have a comprehensive theory to analyze such systems and to provide the basis for their design. This is the object of alternating-current circuit theory. It differs from direct-current theory in that there are circuit elements other than resistors (particularly inductors and capacitors) which respond to (impede) the flow of current: these are known generally as impedances. Also, it is possible to couple energy from one circuit to another by means of mutual inductances or transformers without the need for a linking conductor. Whereas in the dc case the sizes of voltages, currents, and resistors can all be described by single numbers, in the ac case pairs of numbers are needed to describe amplitude and phase or, for an impedance, amplitude ratio and phase shift. Mathematically, the properties and interrelations of pairs of numbers can be dealt with succinctly by using complex algebra. Alternating-current circuit theory is treated using simple time-domain analysis a
a = A 0 sin ωt
T 2
T 2
T
t
A0 (a)
a
sin θ
T
T 2
A0 a = A 0 sin ( ωt + θ)
θ ω
T 2
T
A0 (b) Fig. 2. Sinusoidal waveforms. (a) Simple sinusoid without phase shift (b) Sinusoid with phase shift.
tude) of a and ω is known as the angular or radian frequency. The period T of the function is the time after which the function repeats itself. During one period the function takes on all its possible values and is said to go through one cycle. The frequency f of the sinusoid as usually defined (for example, 60 Hz) is the number of cycles that occur in 1 s, so f = 1/T. Since sin (x + 2π) = sin x, ωT = 2π, and hence ω = 2πf. It is necessary to consider waveforms which do not cross the zero level at t = 0 but which include a phase shift. Figure 2b shows the waveform of Eq. (3), which waveform is identical in shape with a(t) = A0 sin (ωt + θ)
(3)
the waveform described by Eq. (2) but shifted in time so that it crosses the zero level at a time θ/ω before t = 0. This is referred to as a phase advance; the corresponding waveform with θ negative would illustrate a phase delay. Since sin (x + π/2) = cos x, Eq. (3) could equally well be written as Eq. (4), where the phase angle is π π a = A0 sin ωt + + θ − 2 2 = A0 cos (ωt + ψ) (4) now ψ = θ − π /2 radians. Henceforth, either of the forms in Eqs. (3) and (4) will be used, as convenient, to represent an alternating quantity. An alternating voltage is written as Eq. (5), and an alternating current as Eq. (6), where υ and i represent the instanta-
A0
T
(2)
t
v = V0 sin (ωt + θ) or v = V0 cos (ωt + ψ)
(5)
i = I0 sin (ωt + θ) or i = I0 cos (ωt + ψ)
(6)
neous values, and V0 and I0 the maximum values or amplitudes, of voltage and current respectively. See ALTERNATING CURRENT. Phasor diagrams. It is helpful to visualize an alternating waveform as the projection of a rotating vector. This may be considered as corresponding in some sense to the generation of an alternating current by the rotation of a coil in a magnetic field. In Fig. 3a, OR is a vector of length A0 rotating in the anticlockwise direction about the origin of the x-y plane with angular velocity ω radians per second. Then the projection of the vector OR on the x axis is OX = A0 cos ωt, and the projection on the y axis is OY = A0 sin ωt. The rotating vector can therefore be thought of as generating a sinusoidal waveform. If, instead of the vector OR, the vector OR in Fig. 3b is considered, displaced from OR by an angle φ and rotating with it, then OX = cos (ωt + φ) and
509
510
Alternating-current circuit theory y ω
ω
V0
y
R
Y A0
ωt
x
X
0
I0
φ
x
0 (a) (c)
y ω
y Y' A0
φ
ωt 0
ω
R'
X'
V0
R
x
V2
φ
V1
φ ωt 0 (b)
X1
X2
X' 2
x
(d)
Fig. 3. Phasors. (a) Simple phasor. (b) Phasor with phase shift. (c) Voltage and current phasors. (d) Phasor for the sum of two voltages.
OY = sin (ωt + φ), so that rotation of the phasor (the name given to the vector in this type of diagram) corresponds to a phase shift in the waveform. In a linear circuit, the voltage and current waveforms have related amplitudes and often a phase shift between them, but both are of the same frequency. They can thus be represented, as in Fig. 3c, as the projections of phasors of amplitudes V0 and I0 with an angle φ between them, rotating together in the x − y plane. Circuit theory enables the ratio V0/I0 phase shift φ to be calculated from the circuit parameters. If there are two voltages (or two currents) of different amplitudes and with a phase shift between them, say υ 1 = V1 cos ωt and υ 2 = V2 cos (ωt + φ), represented as the projection on the x-axis of two phasors of amplitudes V1 and V2, it is easy to see from Fig. 3d that their sum can be represented as the projection of the resultant phasor V0, the vector sum of V1 and V2. For OX1 = υ 1, X1 X2 = OX2 = υ 2, and therefore OX2 = υ 1 + υ 2. Complex representation. While a phasor diagram helps to visualize the relationship between current and voltage in an ac circuit, it does little to help with the mathematical calculations. These are greatly simplified if complex algebra is used to describe the current, voltage, and circuit parameters. Only the formulas of complex algebra essential for circuit theory are reviewed here. See COMPLEX NUMBERS AND COMPLEX VARIABLES.
A complex number z is defined by a pair of ordinary numbers. It can be written as in Eq. (7), where z = x + jy
(7)
x and y are real numbers and j = −1. (Mathematicians often use i for this quantity, but in electrical work the symbol i is reserved for electric current.) Here x is called the real part of z, written Re (z), and y the imaginary part, written Im (z). If x and y are regarded as coordinates in a (complex) plane, z defines a point in the plane, as in Fig. 4a. The distance r from the origin O to z is known as the modulus of z, written |z|; and the angle θ between the line from O to z and the real (x) axis is known as the argument of z. These quantities are given by Eqs. (8) and (9), and x and y are expressed in terms of them by Eqs. (10). |z| = r = x2 + y2 (8) θ = tan−1 y/x x = r cos θ y = r sin θ
(9) (10)
De Moivre’s theorem states that Eq. (11) is valid, e jθ = cos θ + jsin θ
(11)
so an alternative way of writing z is Eq. (12),
Alternating-current circuit theory and in giving numerical values it is often convenient to write Eq. (13). z = re jθ
(12)
θ
(13)
z=
r
z1 + z2
z2 y2
The complex conjugate of z, z ∗ is defined by Eq. (14). This is illustrated in Fig. 4b. z∗ = x − jy = re−jθ
(14) z1
Multiplication by j is equivalent to rotation through π/2 radians, or 90◦, as Fig. 4c illustrates, for
y1 x1
x2
imaginary axis
z r θ
real axis
x
0
z1 − z2
y −z 2
Fig. 5. Addition and subtraction of complex numbers.
Eq. (15) is valid since j2 = −1, and the real and imagijz = j(x + jy) = −y + jx (a)
z r y θ θ
x y z*
(b)
jz
x r
nary parts of the number interchange as shown. Similarly, multiplication by −j corresponds to rotation through 90◦ in the opposite (clockwise) direction. From this it is apparent that j = e jπ/2. In circuit analysis it is necessary to combine complex numbers in various ways. The rules of complex algebra for combining the numbers z1 = x1 + j y1 and z2 = z2 + j y2 are as follows. The rule for addition is given by Eq. (16). This can z1 + z2 = (x1 + x2 ) + j( y1 + y2 )
θ
z r θ x
x
y
y
(17)
Similarly, z1 − z2 can be represented as the vector sum of the two vectors Oz1 and O(−z2), as in Fig. 5. The rule for multiplication is given by Eq. (18), and that for division by Eq. (19). The process z1 z2 = (x1 + jy1 )(x2 + jy2 ) = (x1 x2 − y1 y2 ) + j(x1 y2 + x2 y1 ) = r1 r2 e j(θ1 +θ2 )
r
(16)
be interpreted geometrically by saying that z1 + z2 is represented by the vector sum of the two vectors Oz1 and Oz2, as shown in Fig. 5. The rule for subtraction is given by Eq. (17). z1 − z2 = (x1 − x2 ) + j( y1 − y2 )
y
(15)
(18)
z1 x1 + jy1 = z2 x2 + jy2 − jz
(c) Fig. 4. Complex numbers. (a) The number z = x + j y = rejθ. (b) Its complex conjugate z∗ = x − jy = rejθ. (c) Multipli◦ cation of z by ± j, equivalent to rotation through 90 .
=
(x1 + jy1 )(x2 − jy2 ) (x2 + jy2 )(x2 − jy2 )
=
x 2 y1 − x 1 y2 x1 x2 + y1 y2 +j x22 + y22 x22 + y22
= (r1 /r2 )e j(θ1 −θ2 )
(19)
511
512
Alternating-current circuit theory i
illustrated in Eq. (19) for converting the quotient of two complex numbers into a number of form a + jb by multiplying both numerator and denominator by the complex conjugate of the denominator is known as rationalization. A comparison of Figs. 3a and 4a shows that alternating waveforms can be represented by the real or imaginary parts of complex numbers of argument ωt, as in Eqs. (20). A0 cos ωt = A0 Re e jωt A0 sin ωt = A0 Im e jωt
R
V
(a)
ω
(20) V0
A sinusoid of arbitrary phase may be written as in Eq. (21), where B is given by Eq. (22) [by De
(v )
A0 cos (ωt + θ ) = 1 /2 A0 e j(ωt+θ) + 1 /2 A0 e−j(ωt+θ ) = Be
jωt
∗ −jωt
+B e
B = 1 /2 A0 e jθ
(21)
i = Ie jωt = I0 e jθ e jωt
(i ) ωt
(22)
Moivre’s theorem]. B contains all the information needed in circuit analysis, since A0 and θ can be calculated from B, and the frequency is given by ω. It is therefore possible to formulate a circuit analysis in which the driving waveforms are of the form given by Eqs. (23), where V and I are complex quantities. v = V e jωt = V0 e jθ e jωt
v0 R
(23)
The circuit parameters are also in complex form, and from these a response in similar form can be calculated which can be directly interpreted as a sinusoidal waveform. Circuit elements. With constant currents, the only circuit element of significance is resistance. Voltage, current, and resistance are related by Ohm’s law: circuit theory is needed to calculate the effective resistance of various series and parallel combinations of individual resistors. With alternating currents, inductance and capacitance have also to be taken into account. In these elements the rate of change of current or voltage is important. The relation between alternating current and voltage in the various types of circuit elements will be discussed. For purposes of theory, these are regarded as ideal elements of one kind only. In practice it is impossible to make an inductor, for example, without some residual resistance, but this will be ignored in the discussion below. Resistance. As in the dc case, the instantaneous voltage υ across a resistance R and the current i through it are related by Ohm’s law, Eq. (24). If υ is given by Eq. (25), then i is given by Eq. (26). Current and
(b) Fig. 6. Resistance. (a) Circuit diagram of resistor and source of alternating voltage. (b) Phasor diagram of voltage and current through a resistance.
Self-inductance. An inductor typically takes the form of a coil of wire, with or without a core of highpermeability magnetic material depending on the value of inductance required and the frequency of operation. The essential property of a single, or self-, inductance (so called to distinguish it from mutual inductance involving two coils, discussed below) is that the voltage across its terminals is related to the rate of change of current through it by Eq. (27). The
v=L
di dt
(27)
inductance L is measured in henrys when the voltage and current are in volts and amperes. If the current is alternating and given by Eq. (28), the voltage is then given by Eq. (29). Thus the ratio i = I0 cos ωt
(28)
v = LI0 (−ω sin ωt) = ωLI0 cos ωt +
π 2
(29)
of the amplitudes of voltage and current is ωL and the voltage waveform is in advance of the current by a phase shift of π /2, or 90◦ (Fig. 7). If the current is instead represented by a complex quantity as in Eq. (30), then Eq. (31) is obtained.
v = iR
(24)
v = V0 cos ωt V0 i= cos ωt R
(25)
i = Ie jωt
(30)
(26)
v = jωLIe jωt
(31)
voltage are in phase and can be represented in a phasor diagram as in Fig. 6. See ELECTRICAL RESISTANCE; RESISTOR.
This represents the same result, since multiplying by j is equivalent to a positive (anticlockwise) rotation through a right angle. See INDUCTANCE; INDUCTOR.
Alternating-current circuit theory the first coil by Eq. (37), where M is the mutual inductance, measured in henrys.
i
v2 = M
L
V
di1 dt
(37)
If i1 is given by Eq. (38), then υ 2 is given by Eq. (39), and the voltage is 90◦ out of phase with the primary current. (a) ω
i1 = I1 e jωt
(38)
v2 = jωLI1 e jωt
(39)
(v)
For further discussion of mutual inductance. See COUPLED CIRCUITS.
V 0 = ω LI 0 I0
(i )
ωt
(b) Fig. 7. Inductance. (a) Circuit diagram of inductor and source of alternating voltage. (b) Phasor diagram of voltage and current through an inductance.
Capacitance. An inductor stores energy in the magnetic field associated with an electric current. A capacitor stores energy in the electric field produced by the separation of charge in a dielectric medium. Typically a capacitor consists of two conducting plates separated by a dielectric. The current flowing in is related to the rate of change of voltage across the capacitor by Eq. (32), where C is the capacitance in farads.
i=C
Circuit analysis. The main problem of conventional ac circuit theory is the calculation of the relationships between currents and voltages in circuits containing various combinations of the circuit elements described above. In order to analyze any physical electric circuit, it is necessary to construct a mathematical model. Such a model is an idealization of the real circuit, but the behavior of the model may be made to approximate that of the real circuit sufficiently closely so that analysis of the model yields results which may be applied with confidence to the real circuit. In the vast majority of cases the circuits to be analyzed are linear. In a linear circuit the most important property is that if the amplitude of the source voltages or currents is changed, then the amplitudes of all other currents and voltages are changed in the same proportion. In the case of sinusoidal currents and voltages every current and voltage in the circuit has i
V
dv dt
C
(32)
If the voltage is given by Eq. (33), then the current is given by Eq. (34). So in this case the ratio of voltv = V0 cos ωt
(a)
(33)
i = CV0 (−ω sin ωt) = ωCV0 cos ωt +
π 2
(34)
(i ) I0
age to current amplitudes is 1/ωC, and the voltage waveform lags behind the current by π/2 (Fig. 8). In complex form, if the voltage is given by Eq. (35), the current is given by Eq. (36), representing the v = V e jωt
(35)
i = jωCV e jωt
(36)
same result. See CAPACITANCE; CAPACITOR. Mutual inductance. If two coils are linked by a common magnetic flux, the voltage υ 2 across the second coil is related to the rate of change of current i1 in
ωt
I V 0 = ω0c
ω
(v ) (b) Fig. 8. Capacitance. (a) Circuit diagram of capacitor and source of alternating voltage. (b) Phasor diagram of voltage and current through a capacitance.
513
514
Alternating-current circuit theory +
+
When ωt + θ = 0, Eq. (41) reduces to Eq. (42) or Eq. (43).
i
R sin (β − θ) + Lω cos (β − θ) = 0 e
i
tan (β − θ) = − −
− (a)
(b)
(42)
R
e
(c)
Fig. 9. Sources and signs. (a) Voltage source. (b) Current source. (c) Sign convention.
the same frequency, ω. In the rest of this article only linear circuits are discussed. In an electric circuit the elements are connected together in some given fashion. Sources (or generators) of external current or voltage are applied to the circuit, and the problem is to calculate the resulting current and voltage associated with each circuit element. Figure 9a and b illustrates the symbols used for voltage and current sources respectively. Although voltage and current change direction every half-cycle, it is necessary to have a sign convention to relate voltage and direction of current flow, compatible of course with that used in dc circuits. This is illustrated in Fig. 9c which is to be understood in the sense that, during the half-cycle in which the upper terminal of the source in the diagram carries the positive voltage, current flows through the resistor in the direction shown, from top to bottom of the diagram. In the next half-cycle both voltage and current reverse. The power dissipated in the resistor is υι, so the convention ensures that the power is positive in each half-cycle. See GENERATOR. The analysis of any circuit is based on Kirchhoff’s laws, which state that: 1. The sum of the instantaneous currents entering the junction of two or more circuit elements is zero. 2. The sum of the instantaneous voltages around any closed loop formed by two or more circuit elements is zero. In both cases, proper account must be taken of the convention regarding positive and negative quantities. See KIRCHHOFF’S LAWS OF ELECTRIC CIRCUITS. Analysis in the time domain. Consider the simple circuit shown in Fig. 10, which consists of a voltage source, a resistor, and an inductor connected in series. (Components are in series when the same current flows through each.) Kirchhoff’s law indicates that the sum of the voltages around the loop is zero; that is, Eq. (40) holds. Now the voltage source e vR + vL − e = 0
I0 [R cos (β − θ) − Lω sin (β − θ)] = E0
L2 ω 2 R2 I0 √ +√ = E0 R2 + L 2 ω 2 R2 + L 2 ω 2 E0 I0 = √ R2 + L2 ω2
(45) (46) (47)
Hence the voltages across the resistor and the inductor are given by Eqs. (48) and (49). The pha
E0 R ωL · sin ωt + θ − tan−1 (48) vR = √ R R2 + ω2 L2 vL = √
E0 ωL R2 + ω2 L2
ωL π · sin ωt + θ − tan−1 + R 2
(49)
sor diagram for this circuit is given in Fig. 3b. The voltages across the resistor and the inductor are π /2 R
+
+
VR
i −
+
vL
e
L
− − (a)
e = E0 sin (ωt + θ)
VL
E0 VR φ
is E0 sin (ωt + θ ), and the (unknown) current is assumed to be I0 sin (ωt + β), so that Eq. (41) is valid, and must be solved to find I0 and β.
− E0 sin (ωt + θ) = 0 (41)
(43)
When ωt + θ = π/2, Eq. (41) reduces to Eq. (44) or (45). Substituting from Eq. (43) for sin (β − θ ) and cos (β − θ) yields Eqs. (46) or (47). π RI0 sin +β −θ 2 π + LωI0 cos + β − θ = E0 (44) 2
(40)
RI0 sin (ωt + β) + LωI0 cos (ωt + β)
Lω R
(b)
I0 VR VL φ
= = = =
I
ωt + θ
E 0 / (R2 + ω2L2) I 0R I 0 ωL tan −1 (ωLIR )
Fig. 10. Circuit with resistance and inductance. (a) Circuit diagram. (b) Phasor diagram.
ω
Alternating-current circuit theory radians out of phase, and the voltage across the circuit leads the current through it by the phase angle tan−1 (ωL/R). The analysis of any other circuit proceeds in the same way. The necessary equations are written by using Kirchhoff’s laws and the relationships of Eqs. (24), (29), and (34) as appropriate. The solutions for the unknown currents or voltages are sinusoids with the same frequency as the applied voltages and currents but unknown amplitudes and phase angles. By evaluating the equations at particular values of time, as in the above examples, the unknown amplitudes and phase angles are then determined. The phase angles are found as differences relative to the phase angle of the applied generator. In the above example, β − θ is found, but not β as such. This is because one phase angle (here θ ) must be taken as a reference angle, and the only important quantity is the value of other phase angles relative to this. Another way of looking at this is to note that the choice of origin on the time axis is quite arbitrary since all voltages and currents extend from t = − ∞ to t = ∞, which is an idealization of the real situation where the generator must start at some definite time. The choice of the reference phase angle corresponds to a particular choice of time origin. While this method can, in principle, be used to analyze any circuit, it becomes extremely cumbersome for a circuit of even moderate complexity. Simplification of the analysis by the use of complex algebra will be discussed, but first the dissipation of power in an ac circuit must be considered. Power. Consider the current and voltage associated with an arbitrary circuit (Fig. 11). The instantaneous power p is equal to υi. Let the voltage and current be given by Eqs. (50) and (51). Then the instantaneous v = V0 sin (ωt + θ)
(50)
i = I0 sin (ωt + β)
(51)
power is given by Eq. (52), which is the sum of a p = V0 I0 sin (ωt + θ ) sin (ωt + θ ) = 1 /2 V0 I0 [cos (2ωt + θ + β) + cos (β − θ)] (52) sinusoid at radian frequency 2ω and a constant independent of t. The instantaneous power is generally of much less interest than the average power P. The average value of a sinusoid is zero so that the aver-
I
+
v
− Fig. 11. Generalized circuit.
circuit
√ age power is given by Eq. (53), where Vˆ = V/ 2 and P = 1 /2 V0 I0 cos (β − θ) = Vˆ Iˆ cos (β − θ) (53) √ Iˆ = I/ 2 are called the root-mean-square (rms) values of υ and i respectively, and cos (β − θ) is called the power factor of the circuit. Thus the power depends not only on the amplitudes of the current and voltage but also on the phase angle between them. This is illustrated by the cases of the resistor, given by Eqs. (54); the inductor, given by Eqs. (55); and the capacitor, given by Eqs. (56). For both the inductor and the capacitor the average power is zero. V0 = RI0 , β − θ = 0 resistor (54) 2 2 P = Vˆ Iˆ = RIˆ = Vˆ /R V0 = LωI0 , β − θ = − π2 inductor (55) P=0 I0 = CωV0 , β − θ = π2 capacitor (56) P=0 It follows therefore that in an ac circuit, power is dissipated only in resistive elements, not in inductors or capacitors. The voltage across these latter elements is 90◦ out of phase with the current through them. The total voltage across a series circuit is made up of the vector sum of the phasors representing the individual element voltages. It is common to speak of in-phase components of the voltage (having the same phase as the current) and out-of-phase or quadrature components with a 90◦ phase difference. Because no power is dissipated in these components, the current through an inductor or capacitor is sometimes referred to as a wattless current. In electric power systems the values quoted for current and voltage are usually the rms values since the power available is related directly to these quantities. However, it must be borne in mind that they must be multipled by 2 to obtain the actual amplitudes of current and voltage. Analysis using complex impedances. The method of time domain analysis described above is quite general, and any circuit may be solved with those techniques. However, the method is cumbersome and does not give much insight into circuit behavior. Further, and perhaps more important, there are a wide variety of circuit theorems developed for use in dc circuit analysis which can both simplify analysis and provide considerable insight, but which cannot be used in ac circuits if one is restricted to the methods of time domain analysis. These considerations led to the development of a method of analysis based on complex notation which makes these techniques available in ac circuit analysis. In an ac circuit every current and voltage is a sinusoid of the same frequency ω, and each is characterized by two quantities, namely the amplitude and the phase angle relative to some reference angle. Thus two independent qualities are needed to describe each current and voltage. As discussed above, an alternating current or voltage can be represented by
515
516
Alternating-current circuit theory an expression of the form Ie j ωt or Ve j ωt, and this representation can be applied to circuit analysis. In the previous example of Fig. 3, Eq. (40) leads to Eq. (57) or (58), from which the modulus and argument of I are given by Eqs. (59) and (60). Thus if the voltage source is given by Eq. (61), then the current is given by Eq. (62), which is the result previously RIe jωt + jLωIe jωt − Ee jωt = 0
I
+
V
Z
_
(57)
Fig. 12. General impedance.
E I= R + jLω |I| = √
(58) I
E R2
+
(59)
L2 ω 2
Argument I = Argument E − tan−1
Lω =β R
e = E0 sin (ωt + θ)
Lω i = |I| sin ωt + θ − tan−1 R
Z2
Z3
V
(60) (61)
_ (a)
(62)
found. However, now the analysis is very simple by using the complex quantities. In the case of the resistor, inductor, and capacitor and in the example, equations are obtained where every term is multiplied by ejωt, and this is true in general. Thus, this multiplier could just as well be omitted in every case and equations could be written involving only the complex quantities V and I. Once V and I are found, the answer for sinusoidal generators can be found immediately. Thus, for a resistor, V = RI; for an inductor, V = jωLI; and for a capacitor, V = I/j Cω. For the complex voltages and currents corresponding to a frequency ω, an inductor behaves as if it were a resistor of value jωL, and a capacitor behaves as if it were a resistor of value 1/jωC. The problem of analysis is thereby reduced to the analysis of dc circuits, except that all currents and voltages are complex numbers and some “resistors” are real while others are imaginary. See NETWORK THEORY. When these complex currents and voltages are used, the ratio V/I in a circuit is called the (complex) impedance of the circuit (usually written Z; Fig. 12), and the ratio I/V is called the (complex) admittance of the circuit (usually written Y ). The units of impedance are ohms, and of admittance are siemens (mhos). The impedance and admittance are expressed in terms of their real and imaginary parts by Eqs. (63) and (64) and are related to each other by Eq. (65) or (66). Similarly, Eq. (67) is valid.
I +
V
Y1
Y2
Y3
_ (b) Fig. 13. Combination of impedances and admittances. (a) Impedances in series. (b) Admittances in parallel.
R + jX =
G 1 B = 2 −j 2 G + jB G + B2 G + B2
(67)
Impedances, like resistances, can be combined in series and parallel. Figure 13a shows three impedances in series. By definition the voltage across each impedance is ZI, and application of Kirchhoff’s law to the loop yields Eq. (68) or (69), so the equivaZ1 I + Z2 I + Z3 I − V = 0
(68)
V = Z1 + Z2 + Z3 = Zeq I
(69)
lent impedance Zeq is the sum of the impedances in series. By definition the current in each admittance is YV, and application of Kirchhoff’s law yields Eq. (70) or Eq. (71). Thus the reciprocal of the equivalent
V = Z = R + jX I
(63)
Y1 V + Y2 V + Y3 V − I = 0
(70)
I = Y = G + jB V
(64)
1 V = Zeq = I Y1 + Y2 + Y3
(71)
1 Z
(65)
X R 1 −j 2 = 2 R + jX R + X2 R + X2
(66)
Y = G + jB =
Z1
+
impedance Zeq is the sum of the admittances in parallel. The real and imaginary components of impedance and admittance are given special names. For impedance, Z = R + jX, the real part R is resistance
Alternating-current circuit theory in the usual way; the imaginary part X (ωL for an inductor or −1/ωC for a capacitor) is reactance. For admittance, Y = G + jB, the real part G is called conductance, and the imaginary part B, susceptance. In a series circuit having resistance R and reactance X, the conductance is not simply the reciprocal of the resistance but, from Eq. (66), a quantity involving both the resistance and the reactance, namely R/(R2 + X2). Occasionally, it is useful to be able to refer to either impedance or admittance without specifying which in particular is meant. The general term immittance is used to cover both. For obvious reasons it is not possible to specify units for it. See ADMITTANCE; CONDUCTANCE; ELECTRICAL IMPEDANCE; IMMITTANCE; REACTANCE; SUSCEPTANCE. Power dissipated. As discussed above, the power entering a circuit, such as Fig. 11, is given by Eq. (72), P = 1 /2 V0 I0 cos (θ − β)
(72)
where V0 and I0 are the amplitudes of the voltage and current sinusoids, and θ − β is the angular phase difference between them. This power can also be expressed in terms of the complex quantities V and I. If V and I are written as in Eqs. (73) and (74), then Eqs. (75) and (76) are valid. Thus, the power is given by Eq. (77).
∗
V I = V0 I0 e
V = V0 e jθ
(73)
I = I0 e jβ
(74)
j(θ−β)
= V0 I0 [cos (θ − β) + jsin (θ − β)]
(75)
V ∗ I = V0 I0 [cos (θ − β) − jsin (θ − β)] (76) P = 1 /4 (V I ∗ + V ∗ I) = 1 /2 V0 I0 cos (θ − β) (77) If V = ZI, then the power is also given by Eq. (78), and if I = YV, then it is also given by Eq. (79). P = 1 /4 (ZII ∗ + Z ∗ II ∗ ) = 1 /4 I02 (Z + Z ∗ )
(78)
P = 1 /4 (Y V V ∗ + Y ∗ V V ∗ ) = 1 /4 V02 (Y + Y ∗ )
(79)
These equations give alternative ways of calculating the power from the voltage or current and circuit parameters. Resonant circuit. Consider the circuit of Fig. 14a with a voltage generator driving a capacitor in parallel with an inductor and resistance in series. The generator produces a voltage e = E0 sin (ωt + θ ). Suppose it is necessary to calculate the currents in the capacitor, resistor, and inductor, and the power dissipated in the circuit. The circuit can be redrawn as in Fig. 14b Y1 as the admittance of the capacitor and Y2 = 1/Z2, where Z2 is the impedance of the resistor and inductor in series. Now Z2 is given by Eq. (80), so Y2 is given by Eq. (81) and Y1 is given by Eq. (82). The equivalent impedance of the circuit, Zeq, is given by Eq. (83), and Z2 = R + jLω 1 Lω R −j 2 = 2 Y2 = R + jLω R + L2 ω 2 R + L2 ω 2
(80) (81)
+
i
R
e
L
C
− (a)
+
I
E
Y1
Y2
− (b) Fig. 14. Resonant circuit. (a) Diagram showing components. (b) Diagram showing admittances.
Y1 = jCω 1 Zeq = Y1 + Y 2
(82) (83)
the generator voltage is related to the total current by Eq. (84) or Eq. (85) so that the total current is given by Eq. (86). E = Zeq I
(84)
I = E/Zeq
(85)
I = E(Y1 + Y2 )
Lω R + j Cω − (86) =E R2 + L2 ω2 R2 + L2 ω2 The current through the capacitor is Y1E = jωCE, while the current through the series combination of the resistor and inductor is Y2E, given by Eq. (87).
R Lω − j (87) Y2 E = E R2 + L2 ω2 R2 + L2 ω2 The power dissipated in the circuit is found, by using Eq. (79), to be given by Eq. (88). ∗ ) P = 1 /4 E02 (Yeq + Yeq
= 1 /4 E02 (Y1 + Y2 + Y1∗ + Y2∗ ) R = 1 /2 E02 2 R + L2 ω 2
(88)
This circuit exhibits the phenomenon of resonance. Suppose the frequency ω is varied. Then in Eq. (86) the imaginary part becomes zero when Eq. (89) is satisfied, which gives Eq. (90) for ω2. This Lω + L2 ω 2
2 1 R = ω02 − ω2 = LC L Cω =
R2
(89)
(90)
517
518
Alternating-current generator is called the resonant frequency, and at this value of ω the power dissipated is given by Eq. (91). The quanP0 = 1 /2 E02
RC L
(91)
tity Lω0/R is called the quality factor or Q-factor of the circuit and is given by Eq. (92). Since Q is often very large, it is usually approximated by Eq. (93). L Q= −1 (92) R2 C 1 L (93) Q= R C Having found these various complex voltages and currents, the corresponding sinusoidal quantities could be written immediately, but this is rarely done in practice since all the information is already contained in the complex quantities. See RESONANCE (ALTERNATING-CURRENT CIRCUITS). Practical circuit elements. Alternating-current circuit theory is, of course, applicable at higher, radio frequencies so long as the circuit elements can be considered as “lumped,” that is, so long as their dimensions are small compared with a wavelength of the current and its phase is the same at all points in the component. But particularly at the higher frequencies, the idealized representations of circuit elements which were given earlier need to be replaced by representations which more accurately reflect the properties of actual components. An inductor, for example, in the form of a wirewound coil, will have resistance as well as inductance, so it should be shown as a resistance and inductance in series, with impedance R + jωL. For a well-designed coil, R is much smaller than ωL at the frequency of use. There will also be capacitance between the turns of the coil which becomes important at the higher frequencies. As discussed above, this is a resonant circuit, and there is for any coil a frequency at which it is self-resonant; above this frequency its effective impedance is capacitative rather than inductive. See CORE LOSS; SKIN EFFECT (ELECTRICITY). A capacitor will generally suffer some power loss in the dielectric medium, which can be represented by a parallel high resistance, and the connecting leads will have low resistance and inductance which may be significant at very high frequencies. See PERMITTIVITY. Other stray impedances may also have to be taken into account. At high frequencies it is particularly important to remember that even lengths of connecting wire may have significant inductance and capacitance to ground. Microwave circuits. At microwave frequencies the assumption of lumped components whose size is small compared with the wavelength breaks down. The basic principles of ac circuit theory can still be applied, but it is necessary to assume that the properties of resistance, inductance, and capacitance are no longer localized but are distributed throughout the circuits, and to take account of the finite time
of travel of a wave from one part of the circuit to another. The concept of impedance then has to be generalized. See MICROWAVE; TRANSMISSION LINES. Active circuits. In addition to the passive circuits discussed in this article, circuit theory can be extended to cover the cases of amplification and feedback. Again the concept of impedance may have to be extended: it is possible, for example, to produce circuits whose effective impedance is that of a negative capacitor, or which do not correspond to any possible combination of simple circuit elements. See AMPLIFIER; CONTROL SYSTEMS; FEEDBACK CIRCUIT. Nonsinusoidal waveforms. The above analysis of circuits has been entirely in terms of the response to single-frequency sinusoidal waveforms. There are many cases where it is necessary to deal with other waveforms. Two methods are available: 1. When it is necessary to find the response to a repetitive waveform or perhaps a single pulse whose Fourier transform can easily be calculated as F(ω), the impedance or transfer function (circuit property relating input and output voltage or current) can be calculated in the usual way as Z(ω), and then the inverse transform of F(ω) · Z(ω) can be extracted to give the response waveform as a function of time. See FOURIER SERIES AND TRANSFORMS. 2. An alternative method, particularly useful in dealing with transient waveforms, is (for example) to take the Laplace transform of the input waveform, L(s), and to multiply it by Z(s), given by the impedance or transfer function calculated in the usual way but with jω replaced by s. The output is then the inverse Laplace transform of L(s) · Z (s). See LAPLACE TRANSFORM. J. O. Scanlan; A. E. Bailey Bibliography. C. K. Alexander and M. N. O. Sadiku, Fundamentals of Electric Circuits, 3d ed., McGrawHill, 2007; J. W. Nilsson and S. A. Riedel, Electric Circuits w/PSpice, 7th ed., Prentice Hall, 2005.
Alternating-current generator A machine that converts mechanical power into alternating-current electric power. Almost all electric power is produced by alternating-current (ac) generators that are driven by rotating prime movers. Most of the prime movers are steam turbines whose thermal energy comes from either fossil or nuclear fuel. Combustion turbines are often used for the smaller units and in cases where gas or oil is the available fuel. Where water power is available from dams, hydroelectric ac generators are powered by hydraulic turbines. Small sites may also use diesel or gasoline engines to drive the generator, but these units are usually used only for standby generation or to provide electric power in remote areas. See DIESEL ENGINE; GAS TURBINE; HYDRAULIC TURBINE; HYDROELECTRIC GENERATOR; INTERNAL COMBUSTION ENGINE; STEAM ELECTRIC GENERATOR; STEAM TURBINE. Alternating-current generators are used because ac power can easily be stepped up in voltage, by using transformers, for more efficient transmission of power over long distances and in larger amounts.
Alternating-current generator
B = Bmax cos
pθ = Bmax cos θe 2
(1)
+ ec _
stator
air gap
B Bmax
θ
p θ 2
electrical radians
(2)
is the number of poles. This angle is defined so that 360 electrical degrees corresponds to an angle that passes through flux defining a north pole and then a south pole. See ELECTRICAL DEGREE. The total flux linking the coil is given by Eq. (3), where the differential area is given by Eq. (4), and φc = BdA (3) dA = LrDθ =
2Lr dθe p
(4)
L is the coil length, r is the radius of the air gap in the cylindrical geometry of the machine, and θ is the angle defined above. The generator shaft rotates at synchronous speed with velocity given by Eq. (5), where ωe, defined by Eq. (6), is the speed in electriωs =
2πf 2 = ωe radians per second p/2 p ωe = 2πf
(5) (6)
cal radians per second and f is the system frequency. The flux density may be regarded as a traveling wave, given by Eq. (7). Substituting the expressions p B(θ, t) = Bmax cos (θ − ωs t) 2 = Bmax cos (θe − ωe t)
(7)
+π 2
+ρ 2
0
−ρ 2
−π 2
Fig. 1. End view of one coil of an alternating-current generator, linked by the moving airgap flux.
in Eqs. (4) and (7) for dA and B into the integrand of Eq. (3) and evaluating between the limits ±ρ/2 yields Eq. (8). Then by using the identity of Eq. (9), ρ/2 2Lr φc = (8) [Bmax cos (θe − ωe t)] dθe −ρ/2 p
sured around the air gap of the machine in a certain direction. The angle θ e is defined by Eq. (2), where p θe =
rotor
ωs
φ
Similarly, transformers step the voltage down again at the utilization site to safer and more convenient levels. See DIRECT-CURRENT GENERATOR; ELECTRIC POWER SYSTEMS; TRANSFORMER. Principles of operation. Most ac generators are synchronous machines, that is, the rotor is driven at a speed that is exactly related to the rated frequency of the ac network. Generators of this type have a stationary armature with three windings that are displaced at regular intervals around the machine to produce three-phase voltages. These machines also have a field winding that is attached to the rotor. This winding provides magnetic flux that crosses the air gap and links the stator coils to produce a voltage according to Faraday’s law. The field winding is supplied with direct current, usually through slip rings. See ARMATURE; SLIP RINGS; WINDINGS IN ELECTRIC MACHINERY. To understand the action of an ac generator, it is helpful to visualize a rotating magnetic flux-density wave in the air gap of the machine (Fig. 1). This wave links the stator winding, causing each coil to experience an alternating flux. This is the mechanism for inducing an alternating voltage. It is usually assumed that the flux density in the air gap has a sinusoidal distribution, which may be written as Eq. (1), where θ is the angular position mea-
519
cos (x − y) = cos x cos y − sin x sin y
(9)
the integration is readily performed. The flux for one coil is found to be given by Eq. (10), where kp and φ p are defined by Eqs. (11) and (12). φc = kp φp cos ωe t kp = pitch factor = sin φp = flux per pole =
(10) ρ 2
4Bmax Lr p
(11) (12)
The induced electromotive force (emf) for the coil is computed from Faraday’s law, which states that the emf is equal to the rate of change of flux linkages; that is, Eq. (13) is valid where Nc is the number of turns in ec = −
d(Nc φc ) dλc =− = ωe Nc kp φp sin ωe t dt dt
(13)
the coil. It is convenient to write the coil voltage as in Eq. (14), where Ec is the root-mean-square (rms) √ (14) ec = 2Ec sin ωe t value of the coil voltage. The total pitch of the coil (π − ρ) is less than one pole pitch (π). This difference has the effect of reducing harmonics more than it reduces the fundamental component of voltage. This reduction is expressed in terms of the pitch factor. See ELECTROMAGNETIC INDUCTION. The voltage ec is the induced voltage in only one coil. (Fig. 1). Usually, the voltages of a group of coils are added together by connecting the coil ends in
520
Alternating-current generator Ec γ
Ec γ
Ec
E group γ
Ec Fig. 2. Phasor diagram for the total root-mean-square emf of a group of coils, Egroup.
series. The coils in the group are not all in the same slots, however, but are displaced by the slot pitch γ . Therefore, the voltage induced in the individual coils will be out of phase by this angle. This means that the addition of the voltages is not a simple arithmetic addition, but is usually performed as a phasor addition to compute the total rms emf of the group of coils, Egroup. This addition is shown in Fig. 2, where the number of coils, n, in the group is assumed to be four. See ALTERNATING-CURRENT CIRCUIT THEORY. From the geometry of Fig. 2, Eq. (15) may be comEgroup = nEc
sin nγ2 = nEc kd n sin γ2
(15)
puted, where a new constant kd is defined for convenience. Finally, the total phase voltage is composed of p groups in series and is thus given by Eq. (16), Ephase = pEgroup
Fmax
P=
(16)
where p is the number of poles. For steam-turbinedriven generators p is usually 2 or 4. Hydroelectric generators may have a much larger number of poles, depending on the shaft speed. Three-phase generators produce three-phase voltages that are equal in magnitude but displaced in phase by 120 electrical degrees. The phasor diagram for the voltages of a three-phase ac generator is shown in Fig. 3. Armature reaction. Armature reaction is a demagnetizing effect that limits the output of an ac generator. The magnetomotive forces (mmf’s) of the coils combine to produce a rotating mmf that opposes the field flux. Thus, more excitation is required to maintain the flux magnitude as the armature current is increased. The mmf of one coil per pole is a square wave that has a sine-wave fundamental component with a maximum value given by Eq. (17), where Ic is the rms armature current in the √ 4 = 2 N c Ic k p π
currents, the forward components combine and the reverse components cancel. Thus, for three-phase windings, the resultant rotating wave has a maximum of 3/2 the single phase value and is usually referred to as F1, where the subscript 1 indicates the fundamental component of the sine wave. See ARMATURE REACTION. Characteristics. Certain constants define the operation of the ac generator on the system to which it is connected. One important set of machine constants is the reactances of the windings. There are several different reactances that are derived to describe the machine behavior under different conditions. See REACTANCE. The synchronous reactance xd is the unsaturated steady-state reactance, including leakage and armature reaction, and is used to describe steady-state generator performance. The transient reactance xd is useful in determining how the machine responds to sudden changes, after the first few cycles. This reactance is often used to compute the fault current for circuit-breaker opening. The subtransient reactance xd determines the initial fault current, which is a current that the machine windings and connecting bus-work must be able to withstand mechanically. The generator operates in synchronism with the system to which it is connected. As load is increased, the shaft angle moves ahead by the angle δ. The power delivered is given by Eq. (18), where the subEs Eg sin δ xs + xd
(18)
script s stands for system quantities and the subscript g stands for generator quantities. The voltage Eg for steady-state conditions is the rated excitation voltage, and xg is the synchronous reactance. To calculate system swings due to sudden changes, xd can be taken as the transient reactance and Ed as the effective voltage behind that reactance. This
Ec
120°
Ea
120° 120°
(17)
coil. The mmf factor for all coils in one phase is Fmaxkd. To compute the resultant rotating mmf for a polyphase winding, the single alternating wave is divided into two equal components of one-half value each, rotating in opposite directions. For balanced phase
Eb Fig. 3. Phasor diagram of the balanced three-phase generated voltages, Ea, Eb, and Ec, for a synchronous alternating-current generator.
Alternating-current motor voltage changes only slowly but can be corrected for the action of the excitation system, which affects the control of the voltage and plays an important role in the dynamic performance of the machine. The calculation of dynamic swings during a short circuit can be performed step by step or by computer simulations, the latter method being the usual choice for systems containing more than one machine. The computation gives the change in power output of the machine and the resulting acceleration and angle, taking into account the moment of inertia of the turbine-generator combination. Growth in size. The size of ac generators increased rapidly up to the mid-1970s. This increase resulted in considerable economy in terms of the cost per megawatt of machine rating and a concurrent improvement in turbine efficiency. The maximum ratings doubled about every 10 years for many years, which is about the same rate at which the total electrical load in North America was increasing. The largest units now manufactured are rated at about 1800 megavolt-amperes. This rapid growth in unit size was a major challenge to machine designers. Better rotor materials were required to withstand higher centrifugal forces associated with larger rotor diameters. Major changes in the method of cooling were required, resulting first in hydrogen cooling and later in water-cooled stator windings. These refinements were necessary to keep the efficiency of the generator at about 99%. See VOLT-AMPERE. Other generator types. The foregoing discussion has dealt entirely with the class of generators called synchronous machines. These are by far the most common, but are not the only type of ac generators. Induction generators have been used in remote applications where maintenance of the excitation system is a problem. These units are essentially like induction motors, but are driven by a prime mover at speeds slightly above synchronous speed, forcing the unit to generate power due to the reverse slip. Induction generator units draw reactive power from the system and are not as efficient as synchronous generators. See INDUCTION MOTOR. High-frequency single-phase generators have been built as induction alternators, usually with twice as many stator poles (teeth) as rotor poles, and with a constant air-gap flux supplied from a homopolar field coil in the center of the machine, pushing flux into the stator at one end and out at the other. Their effectiveness is lower than that of the synchronous machine because the flux is a pulsating unidirectional field, rather than an alternating field. See ALTERNATING CURRENT; ELECTRIC ROTATING MACHINERY; GENERATOR. Lee A. Kilgore; Paul M. Anderson Bibliography. D. G. Fink and H. W. Beaty (eds.), Standard Handbook for Electrical Engineers, 14th ed., 2000; A. E. Fitzgerald, C. Kingsley, and S. D. Umans, Electric Machinery, 6th ed., 2003; G. McPherson and R. D. Laramore, An Introduction to Electrical Machines and Transformers, 2d ed., 1990.
Alternating-current motor An electrical machine that converts alternatingcurrent (ac) electric energy to mechanical energy. Alternating-current motors are widely used because of the general availability of ac electric power and because they can be readily built with a variety of characteristics and in a large range of sizes, from a few watts to many thousands of kilowatts. They can be broadly classified into three groups—induction motors, synchronous motors, and ac series motors: Induction motor Single-phase Split-phase Capacitor-start Capacitor-run Polyphase Synchronous motor Single-phase Permanent-magnet (PM) Reluctance Hysteresis Polyphase Wound-field Permanent-magnet (PM) Reluctance AC series or universal motor (singlephase) See
ALTERNATING MOTOR.
CURRENT;
DIRECT-CURRENT
Induction motor. The most common type of ac motor, both in total number and in total power, is the induction motor. In larger sizes these machines employ a polyphase stator winding, which creates a rotating magnetic field when supplied with polyphase ac power. The speed of rotation depends upon the frequency of the supply and the number of magnetic poles created by the winding; thus, only a discrete number of speeds are possible with a fixed frequency supply. Currents are induced in the closed coils of the rotor for any rotor speed different from the speed of the rotating field. The difference in speed is called the slip speed, and efficient energy conversion occurs only when the slip speed is small. These machines are, therefore, nearly constant-speed machines when operated from a constant-frequency supply. They are, however, routinely started from zero speed and accelerated through the inefficient high-slip-speed region to reach operating speed. See SLIP (ELECTRICITY). Single-phase induction machines are generally inferior in conversion efficiency and have zero torque at stand-still (zero speed). Auxiliary starting means usually involving a second winding and some technique for creating a time phase difference in the second winding current must be employed. Capacitors are commonly used to create a large phase shift, although extra resistance in the starting winding (in the split-phase motor) can be used if reduced starting
521
522
Alternative fuels for vehicles torque is acceptable. Typically, the second winding is switched off by a rotation-speed-actuated switch, but in some low-starting-torque applications the second winding with a series capacitor is left connected to improve the running condition (in the capacitor-run motor). Single-phase machines are widely used in situations where only single-phase ac power is available, especially for appliances in homes and small commercial installations. See INDUCTION MOTOR. Synchronous motor. In contrast to an induction motor, the rotor of a synchronous motor runs exactly at the rotating field speed and there are no induced rotor currents. Torque is produced by the interaction of the rotating field with a direct-current (dc) field created by injected dc rotor current or permanent magnets, or with a rotor magnetic structure that has an easy direction for magnetization (in the reluctance motor). Since for any frequency of excitation there is only one speed for synchronous torque, synchronous machines have no starting torque unless the frequency is variable. When the motor is used in fixed-frequency applications, an induction-machine winding is also placed on the rotor to allow starting as an induction motor and running as a synchronous motor. See SYNCHRONOUS MOTOR. Large synchronous motors for industrial applications have dc windings on the rotor and always have polyphase windings on the stator. These machines characteristically have high efficiency and a power factor that is adjustable via the rotor dc field current. Both leading and lagging power factors are possible. Smaller polyphase synchronous motors using permanent magnets on the rotor are widely used in self-synchronous, variable-frequency drive systems (brushless dc motors and electronically commutated motors). In these drives the rotor position is detected and used to determine the frequency and phase of the excitation. The advantages include high efficiency, fast response, and very good control characteristics. In smaller sizes, synchronous reluctance motors are used where precise speed control is required, such as in timing motors or where many small motors must be operated at exactly the same speed. Reluctance machines generally have lower efficiency and power factor but are inexpensive. Improvements in design may make the reluctance machine competitive with induction machines in variable-speed applications using variable-frequency excitation. See RELUCTANCE MOTOR. For special, low-power applications requiring precise speed control and very smooth torque, the hysteresis motor is used. These machines use a special permanent-magnet rotor material which is remagnetized each time the motor is started. Starting torque results from the lag between magnetizing force and flux density and is essentially independent of rotor speed. At synchronism the hysteresis motor runs as a permanent-magnet synchronous motor. See HYSTERESIS MOTOR. AC series motor. A dc motor with the armature and field windings in series will run on ac since both magnetic fields reverse when the current reverses. Since these machines run on ac or dc, they
are commonly called universal motors. The speed can be controlled by varying the voltage, and these machines are therefore widely used in small sizes for domestic appliances that require speed control or higher speeds than can be attained with 60-Hz induction motors. See ELECTRIC ROTATING MACHINERY; MOTOR; UNIVERSAL MOTOR; WINDINGS IN ELECTRIC MACHINERY. Donald W. Novotny Bibliography. D. G. Fink and H. W. Beaty (eds.), Standard Handbook for Electrical Engineers, 14th ed., 2000; A. E. Fitzgerald, C. Kingsley, and S. D. Umans, Electric Machinery, 6th ed., 2003; G. R. Slemon, Electric Machines and Drives, 1992.
Alternative fuels for vehicles Conventional fuels such as gasoline and diesel are gradually being replaced by alternative fuels such as gaseous fuels (natural gas and propane), alcohol (methanol and ethanol), and hydrogen. Conventional fuels can also be modified to a reformulated gasoline to help reduce toxic emissions. Technological advances in the automotive industry (such as in fuel cells and hybrid-powered vehicles) are helping to increase the demand for alternative fuels. Two key issues associated with the use of conventional fuels for vehicles prompt interest in alternative fuels: (1) gasoline- and diesel-powered vehicles are considered to be a significant source of air pollution; and (2) conventional fuels are produced from crude oil, a nonrenewable energy resource. The success of alternative fuels in the marketplace will depend on numerous factors, including public understanding and consumer awareness; economics; automobile performance; availability of fuels, vehicles, and distribution and marketing systems; and changes in technology. See AIR POLLUTION. Gaseous fuels. Vehicle emissions from natural gas and propane are expected to be lower and less harmful to the environment than those of conventional gasoline. Because natural gas and propane are less complex hydrocarbons, the levels of volatile organic compounds and ozone emissions should be reduced. Both of these fuels are introduced to the engine as a gas under most operating conditions and require minimal fuel enrichment during warm-up. Leaner burning fuels, they also achieve lower carbon dioxide and carbon monoxide levels than gasoline. However, because they burn at higher temperatures, emissions of nitrogen oxide are higher. An important property of gaseous fuels is their degree of resistance to engine knock. Because of their higher-octane value relative to gasoline, there is less of a tendency for these fuels to knock in spark-ignition engines. To achieve the optimal performance and maximum environmental benefits of natural gas and propane, technological advancements must continue to reduce the costs of dedicated vehicles to be competitive with conventional vehicles, and the necessary fueling infrastructure must be ensured. Natural gas. Natural gas is a colorless, odorless hydrocarbon that is neither carcinogenic nor corrosive.
Alternative fuels for vehicles Once processed at a gas plant, natural gas is composed of about 97% methane (CH4). Being lighter than air, it does not pool on the ground when leaked but dissipates into the atmosphere, reducing the hazard of fire. One drawback as a transportation fuel is its very low energy density relative to gasoline or diesel fuel. To eliminate bulkiness, the gas is compressed and stored in cylindrical tanks. At a pressure of 3000 pounds per square inch (20 megapascals), the volumetric energy density of a typical natural gas cylinder is only 20% that of a similar volume of gasoline. Thus, the driving range of a natural gas vehicle carrying the same volume of fuel as a gasoline-fueled vehicle will be less. Larger vehicles, carrying two or three natural gas tanks, can achieve an adequate driving range for most personal or fleet purposes, but at the cost of space and the additional weight of the cylinders. See METHANE. Natural gas engine technology is fairly well developed, including heavy-duty applications. Vehicle refueling appliances have reduced the constraints associated with the lack of distribution and refueling infrastructure; however, the associated costs in addition to those related to conversion tend to limit natural gas usage to high-mileage vehicles. Significant advances have been made with regard to lighter and stronger fuel cylinders, and there is more variety of natural gas vehicles produced by original equipment manufacturers than for any other alternative fuel. Still required for market competitiveness of natural gas is a refueling infrastructure with widespread access and low-cost fast-fill capabilities. See LIQUEFIED NATURAL GAS (LNG); NATURAL GAS. Propane. Propane (C3H8), a hydrocarbon produced from crude oil and natural gas, has been one of the more successful alternative transportation fuels. A cleaner-burning fuel than gasoline or diesel fuel, it can result in lower maintenance costs, requiring fewer spark plug and oil changes, and less wear on such components as piston rings and bearings. Since propane has a higher octane rating than gasoline, propane-fueled vehicles can use higher engine compression ratios, resulting in more power and better fuel efficiency. Propane has a lower energy density than gasoline by volume but a higher energy density by weight; therefore, the weight of a tank full of propane is similar to that of one filled with gasoline. Compared to ethanol, methanol, or natural gas, propane has a higher energy content by weight and volume, so it can take a vehicle farther on an average-capacity tank than a similar tank of any of the other alternative fuels. Refueling time for propane is similar to conventional fuels. Several factors related to customer convenience and conversion costs have been restrictive to automotive propane market development. The physical characteristics of propane require modifications to conventional vehicles by either the automotive manufacturer at the factory or after-market conversions; either way, the cost of the modifications adds to the price paid by the consumer for a gasoline or diesel vehicle. In order to maintain a liquid state, propane
must be stored under pressures of up to 200 psi (1.3 MPa). Because of this fact and the lower energy density of propane relative to gasoline, propanepowered vehicles must be fitted with storage tanks that impose weight and space penalties in order for the vehicles to travel the same distance as gasoline counterparts. Also, the vapor pressure of propane varies in relation to the atmospheric temperature. At temperatures of −43◦C (−45◦F) and colder, the pressure in a fuel tank is essentially zero, the product becomes liquid, and propane fuel systems simply do not operate. Another inconvenience arises due to the potential hazard of vapor accumulation at ground level; for this reason, some city and municipality bylaws forbid propane vehicles from parking in underground facilities. On the basis of a grams per unit of distance traveled, propane full fuel cycle emissions compare favorably with gasoline, except for nitrogen oxides. With further development of propane vehicles by original equipment manufacturers, the establishment of certification standards for propane conversions, and continued technological improvements in emission control systems, the environmental impact of propane could prove even greater. The potential economic benefits from conversion to propane are greatest in high-fuel-consumption fleets such as taxis, where payback occurs in 1–2 years. Alcohol fuels. The most significant advantage of alcohol fuels over gasoline is their potential to reduce ozone concentrations and to lower levels of carbon monoxide. Another important advantage is their very low emissions of particulates in diesel engine applications. In comparison with hydrocarbon-based fuels, the exhaust emissions from vehicles burning low-level alcohol blends (such as gasohol containing 10% alcohol by volume) contain negligible amounts of aromatics and reduced levels of hydrocarbons and carbon monoxide but higher nitrogen oxide content. Exposure to aldehydes, in particular formaldehyde which is considered carcinogenic, is an important air-pollution concern. The aldehyde fraction of unburned fuel, particularly for methanol, is appreciably greater than for hydrocarbon-based fuels; therefore, catalytic converters are required on methanol vehicles to reduce the level of formaldehyde to those associated with gasoline. See ALCOHOL FUEL. Methanol. Methanol (CH3OH) is a clear, colorless, high-performance liquid fuel that can be used in both spark-ignition and diesel engines. It can be burned as a neat (100% methanol) or near-neat fuel in internal combustion engines. A blend of 85% methanol and 15% gasoline, M85, is the most common form of methanol fuel used in light-duty vehicles. M100, or neat methanol, is used in some heavy-duty trucks and buses. The appeal of methanol as an alternative to gasoline stems from its being a relatively inexpensive clean-burning liquid fuel. A constraint is the fact that it has only half the energy density of gasoline, compensated to some degree by its better thermal efficiency. Safety concerns relate to the fact that
523
524
Alternative fuels for vehicles methanol is toxic and can be absorbed through the skin. Also, neat methanol burns with an invisible flame, making methanol fires without a colorant difficult to detect. Besides being used as a neat or near-neat transportation fuel, methanol has been used to produce methyl tertiary butyl ether (MTBE), an oxygenate which, if blended with gasoline up to 10%, adds one octane number to the fuel. Concerns related to MTBE leaking and contaminating ground water in some areas in the United States resulted in a reexamination of the use of this oxygenate, and the phaseout of its use in California by December 31, 2002. Methanol is also being tested in the transportation industry as a source of the hydrogen in hybrid fuel cell vehicles. See ETHER; METHANOL. Ethanol. Ethanol (C2H5OH) produced from biomass is a renewable fuel source currently being marketed in neat, near-neat, and low-level premium fuel blends. In addition to providing environmental benefits, use of ethanol fuel produced from fermentable starch or sugar crops provides economic advantages to the agriculture sector. However, ethanol production costs are very high relative to gasoline and must be reduced substantially if ethanol is to compete on an economic basis. To date, in the fuel industry ethanol has been used mainly as an additive in low-level gasohol blends. These blends have been successfully used in unmodified gasoline vehicles with warranty coverage provided for automobiles sold in North America. Neat ethanol, generally used in heavy-duty applications, is less common and requires costly engine modifications. Flexible fuel vehicles can operate on a mixture of gasoline and up to 85% ethanol (E85). Ethanol has also been used in the production of ethyl tertiary butyl ether (ETBE), but to date it has not been economically competitive with MTBE. The oxygen content of ethanol fuels used in properly modified vehicles results in increased energy efficiency and engine performance. Extended sparkplug life and lower carbon deposits are also expected because ethanol burns cleaner than gasoline. A significant drawback of ethanol is its much lower energy density than gasoline (approximately 34% less). To travel distances similar to a gasolinepowered vehicle without refueling would require a much bigger fuel tank; this requirement is slightly offset by the greater energy efficiency of ethanol. See ETHYL ALCOHOL. Hydrogen. Hydrogen is the lightest and most abundant element in the universe. It can be produced from a number of feedstocks in a variety of ways. The production method thought to be most environmentally benign is the electrolysis of water, but probably the most common source of hydrogen is the steam reforming of natural gas. Once produced, hydrogen can be stored as a gas, liquid, or solid and distributed as required. Liquid storage is currently the preferred method, but it is very costly. Metal hydride and compressed storage are also being investigated. See METAL HYDRIDES.
Hydrogen-powered vehicles can use internal combustion engines or fuel cells. They can also be hybrid vehicles of various combinations. When hydrogen is used as a gaseous fuel in an internal combustion engine, its very low energy density compared to liquid fuels is a major drawback requiring greater storage space for the vehicle to travel a similar distance to gasoline. Although hybrid vehicles can be more efficient than conventional vehicles and result in lower emissions, the greatest potential to alleviate air-pollution problems is thought to be in the use of hydrogen-powered fuel cell vehicles. Though currently very expensive, fuel cells are more efficient than conventional internal combustion engines. They can operate with a variety of fuels, but the fuel of choice is gaseous hydrogen since it optimizes fuel cell performance and does not require on-board modification. Numerous hydrogen fuel cell vehicle models are being developed and tested with the objective of having such vehicles available for sale to the public in the early years of the twenty-first century. However, for gaseous hydrogen to become a fuel of choice in the transportation industry, numerous technical and economic constraints related to hydrogen production, on-board and site storage, and marketing and distribution systems must be overcome. Safety issues and perception must also be addressed. See FUEL CELL; HYDROGEN. Reformulated fuels. Conventional gasoline and diesel fuels are complex mixtures of many different chemical compounds. Over time these fuels have undergone reformulation to improve their value to the transportation industry; this reformulation ranged from mild, like the backing-out of butane to reduce the volatility, to severe adjustments that substantially change the composition. The U.S. Clean Air Act Amendments (CAAA) have served to increase interest in using regulated changes to motor fuel characteristics as a means of achieving environmental goals. The reformulated gasoline (RFG) program was designed to resolve ground-level ozone problems in urban areas. Under this program, compared to conventional gasoline, the amount of heavy hydrocarbons is limited in reformulated gasoline, and the fuel must include oxygenates and contain fewer olefins, aromatics, and volatile organic compounds. Diesel-powered heavy-duty trucks and buses are said to account for about a quarter of the nitrogen oxides (NOx) that result in smog formation and well over half of the particulates produced by mobile sources. Particulates, unburned fuel particles, can stick to the lungs if inhaled, causing increased susceptibility to respiratory illness. Research is under way to develop diesel fuel that will reduce both the particulate and NOx levels. These fuels will have higher cetane ratings than conventional diesel fuel and lower sulfur and aromatics. Michelle Heath Bibliography. R. L. Bechtold, Alternative Fuels Guidebook: Properties, Storage, Dispensing, and Vehicle Facility Modifications, 1997; R. L. Busby,
Altimeter Hydrogen and Fuel Cells: A Comprehensive Guide, 2005; Society of Automotive Engineers, State of Alternative Fuel Technologies 2000, 2000.
Altimeter Any device which measures the height of an aircraft. The two chief types are the pressure altimeter, which measures the aircraft’s distance above sea level, and the radar altimeter, which measures distance above the ground. Pressure Altimeter A pressure altimeter precisely measures the pressure of the air at the level an aircraft is flying and converts the pressure measurement to an indication of height above sea level according to a standard pressurealtitude relationship. In essence, a pressure altimeter is a highly refined aneroid barometer since it utilizes an evacuated capsule whose movement or force is directly related to the pressure on the outside of the capsule (Fig. 1). Various methods are used to sense the capsule function and cause a display to respond such that the pilot sees the altitude level much as one looks at a watch. Because altitude measured in this manner is also subject to changes in local barometric pressure, altimeters are provided with a barosetting that allows the pilot to compensate for these weather changes, the sea-level air pressure to which the altimeter is adjusted appearing in a window of the dial. Flights below 18,000 ft (5486 m) must constantly contact the nearest traffic center to keep the altimeters so updated. Flights above 18,000 ft and over international waters utilize a constant altimeter setting of
aneroid capsule
29.92 in. Hg, or 1013.2 millibars (101.32 kilopascals), so that all high-flying aircraft have the same reference and will be interrelated, providing an extra margin of safety. See AIR NAVIGATION. Electromechanical computers are provided that correct the sensed static pressure for the static systems defect error, which is an inherent aircraft error due to the effect of high-speed flying on the air mass immediately around the aircraft fuselage in the vicinity of the static ports. The corrected altitude information is usually transmitted to electromechanical indicators for the pilot’s display. Much work has been done to improve the readability of the display, so that pilots may make rapid readings with a low probability of reading error. The present accuracy of altimeters is on the order of 0.5% of the indicated altitude for all mechanical indicators and 0.2% of the indicated altitude for computer systems. See BAROMETER. James W. Angus Radar Altimeter A radar altimeter is a low-power radar that measures the distance of an aircraft (or other aerospace vehicle) above the ground. It differs from a pressure altimeter, which measures distance above sea level when properly adjusted for atmospheric pressure variations. Radar altimeters are often used in aircraft during bad-weather landings. They are an essential part of many blind-landing and automatic navigation systems and are used over mountains to indicate terrain clearance. Special types are used in surveying for quick determination of profiles. Radar altimeters are used in bombs, missiles, and shells as proximity fuses to cause detonation or to initiate other functions at set altitudes. Radar altimeters have been used on
100-ft pointer
barosetting scale
static pressure inlet
1000-ft pointer
10,000-ft pointer
window
graduated scale barosetting knob
Fig. 1. Schematic arrangement of sensitive pressure altimeter. 1 ft = 0.3 m.
525
h distance from radar altimeter
r
h distance from radar altimeter
r
mean return power per unit distance
(a)
(b)
typical single pulse power
various spacecraft, starting with Skylab in 1973, to measure the shape of the geoid and heights of waves and tides over the oceans. Notable spacecraft altimeters over the Earth were those aboard the Seasat and Geosat American satellites, the American-French TOPEX/POSEIDON, and the European ERS 1. Other spacecraft altimeters provide topographic information on other planets, particularly Venus. See AUTOMATIC LANDING SYSTEM; GROUND PROXIMITY WARNING SYSTEM. Principle of operation. Like other radar devices, the altimeter measures distance by determining the time required for a radio wave to travel to and from a target, in this case the Earth’s surface. If the Earth were a perfectly flat horizontal plane or smooth sphere, the signal would come only from the closest point (Fig. 2a), and would be a true measure of altitude. Actually, the Earth is not smooth, and energy is scattered back to the radar from all parts of the surface illuminated by the transmitter (Fig. 2b). For the radar to measure distance to the ground accurately, it must distinguish between the energy from points near the vertical and that from more distant points. Radar altimeters can seldom have highly directive antennas, for they must be able to function regardless of the aircraft attitude. Thus the antenna cannot discriminate against off-vertical signals. The wide beam width means the antenna gain is small, so the power required is much greater than for a comparable directive radar working against a large, close target. Spaceborne altimeters used to measure distances to the ocean and ice sheets, however, are mounted on stable platforms that allow them to use very narrow antenna beams. See ANTENNA (ELECTROMAGNETISM). Most radio altimeters use either pulse or frequency modulation, the former being more popular for high altitudes, and the latter for low altitudes. See FREQUENCY MODULATION; PULSE MODULATOR. Pulse altimeters. In a typical pulse altimeter (Fig. 3a) the radio-frequency carrier is modulated with short pulses (under 0.25 microsecond). The short pulse permits measurements, even at low altitudes, of the time delay between the leading edge of the transmitted pulse and that of the pulse returned from the ground. Early pulse altimeters displayed the received signal on a cathode-ray tube with circular sweep, allowing the pilot to determine the leadingedge position of the echo signal. Modern pulse altimeters use a tracking gate system. One gate is kept close to the leading edge by a servo system that adjusts the position of the gate to the optimum delay point. A simple single-gate system can be used (Fig. 3a), but most pulse altimeters use two or three gates to achieve better distance measurement in the presence of noise and fading. The time delay is displayed for the pilot as a range-dial reading or digital readout. The delay is usually converted to a digital signal that is delivered to the autopilot and flightcontrol computer. The return power for a pulse altimeter (the mean of many pulses; Fig. 2c) is a combination of scattered returns from the entire illuminated area and specular return from the nearest point to which the
return power per unit distance
Altimeter
mean return pulse gate pulse
t
2h/c (c) power at different frequencies
526
(d)
time from start of transmitted pulse
envelope filter pass band
(2h/c)(dfc /dt) difference between received and transmitted frequencies
fd
Fig. 2. Radar altimeter principles. (a) Ideal return power distribution, which would be observed if the signal came only from the closest point to the radar. (b) Typical distribution of average returned power as a function of distance from the radar. (c) Pulse shapes in a wide-beam, narrow-pulse altimeter with a single tracking gate. (d) Typical FM altimeter difference-frequency spectrum. Spacing between lines is the sweep repetition frequency.
beam is perpendicular. The specular contribution is important only for very smooth water, pavements, or desert playas. Usually the maximum value of the return power varies with height h as 1/h2 to 1/h3. Interference between different scattering centers causes fading, so the pulses must be averaged over a considerable time to get a precise indication. See INTERFERENCE OF WAVES. FM altimeters. In a frequency-modulated (FM) altimeter (Fig. 3b), the frequency of a continuous carrier is swept in some manner, usually to give a triangular frequency-time curve. The difference in frequency between that received from the ground (but transmitted earlier) and that being transmitted is
Altimeter antenna pulse source
transmitter
transmitreceive switch
variabledelay gate generator
coincidence circuits
superheterodyne receiver
delayadjusting servo
indicator
(a)
indicator
antenna 1 sweepadjusting servo
sweep generator (adjustable deviation or rate)
frequencymodulated oscillator antenna 2
filter
high-gain amplifier
balanced mixer
(b) Fig. 3. Block diagrams of typical radar altimeters. (a) Automatic pulse altimeter. Control, tracking, and output circuits may all be digital. (b) Servo-adjusted-sweep FM altimeter.
a measure of the time delay (Fig. 2d). This difference frequency is given by the equation below, where r fd =
d fc (2r/c) dt
is the distance from the radar, c is the speed of light, and fc is the carrier frequency. Some radars with fixed dfc/dt (lacking a servo and fixed filter) use electronic frequency meters calibrated in range as indicators. Because the return comes from many ranges besides the altitude, an unsophisticated radar of this type may read an effective range frequency more than 10% high in some cases. Other types (Fig. 3b) use a relatively narrow filter and keep the minimum difference frequency (altitude frequency) centered in this filter by a servo which adjusts the sweep rate dfc/dt. The sweep rate is indicated, as range, on a dial and also digitized for supply to autopilot and flight computer. The frequency spectrum (Fig. 2d) is composed of individual lines, spaced by the sweep frequency. Suitable signal processing permits recovery of the envelope of the spectrum, which is analogous to the mean return pulse of a pulse altimeter (Fig. 2c). Because of problems with short pulses at short ranges,
FM altimeters are normally used where altitude measurements must be below about 100 ft (30 m). Other radio devices may also serve as altimeters. The proximity fuse is a crude device for measuring altitude by field strength alone. Various correlation devices have been proposed, some of which employ noise modulation of either amplitude or frequency. Spaceborne altimeters. Pulse altimeters are used on spacecraft to observe properties of the ocean and continental ice sheets (Greenland and Antarctica). For these applications, the exact location of the spacecraft must be determined by accurate tracking from land. The distance from altimeter to surface then becomes a measure of the variation of the surface itself. The altimeters must achieve great accuracy. To do so, they transmit pulses with very large bandwidths, thereby achieving resolutions of only a few feet. The actual distance measurement accuracy can be better than 5 cm (2 in.), provided suitable corrections are made for slight variations in the refractive index of (and therefore velocity of propagation through) the atmosphere. TOPEX/POSEIDON uses two frequencies to aid in correcting for the minute effects of the ionosphere. To achieve this great precision without the use of extremely short pulses that would need excessive peak power, pulse-compression techniques are employed. Normally these methods use a binary phasecoded pulse. These spaceborne systems have allowed major improvements in knowledge of the shape of the geoid over the oceans, where surveying methods that work with great precision over land perform poorly or cannot be used at all. Moreover, with the precision now available, oceanographers can locate and track changes in the major current systems in the ocean by the slight bulge associated with the current flow. Glaciologists use the altimeter data taken over the continental ice sheets to detect growth and decay of the ice, which can be an indication of global warming or cooling. Since the ice sheets contain most of the fresh water on Earth, significant melting resulting from global warming could result in major catastrophes along coastlines, making monitoring of glaciers an important concern. See GEODESY; GLACIOLOGY; OCEAN CIRCULATION. The same altimeters also can measure wave height at points along the ocean surface beneath the satellite (the subsatellite track). When ocean waves are high, the leading edge of the average pulse (Fig. 2c) is stretched out, since the signal returns first from a relatively small area near the crests of the waves and then builds up as the entire surface is illuminated. The altimeter samples the pulse at close intervals (typically about 10 nanoseconds), and the averaged samples are used to determine the shape of the received pulse. This can be compared with theoretical shapes to achieve good estimates of the root-meansquare wave height beneath the spacecraft. Altimeters can also be used as scatterometers to estimate the wind speed beneath the spacecraft. This is possible because the received signal decreases as the wind increases. Since meteorologists require wind
527
528
Altitudinal vegetation zones measurements over wide areas, this subsatellite wind measurement is less useful than that from spaceborne radar scatterometers looking for hundreds of miles to the side of the spacecraft. Moreover, the wind sensitivity of radar scatterometers beyond an angle of 20◦ from the vertical is much greater than that at vertical incidence. Nevertheless, the altimeters have proved useful during intervals when no spaceborne scatterometer was active. See REMOTE Richard K. Moore SENSING. Bibliography. D. E. Barrick, Analysis and interpretation of altimeter sea echo, Adv. Geophys., 37:60–98, 1985; B. C. Douglas and R. E. Cheney, Geosat: Beginning of a new era in satellite oceanography, J. Geophys. Res., 95:2833–2836, 1990; Seasat Special Issue II: Scientific Results, J. Geophys. Res., 88(C3):1531– 1745, 1983; M. I. Skolnik, Introduction to Radar Systems, 2d ed., 1980; G. I. Sonnenberg, Radar and Electronics Navigation, 6th ed., 1988; F. T. Ulaby, R. K. Moore, and A. K. Fung, Microwave Remote Sensing, vol. 2, 1982.
Altitudinal vegetation zones Intergrading regions on mountain slopes characterized by specific plant life forms or species composition, and determined by complex environmental gradients. Along an altitudinal transect of a mountain, there are sequential changes in the physiognomy (growth form) of the plants and in the species composition of the communities. This sequential zonation of mountain vegetation has been recognized for centuries. Vertical zonation was fully developed as an ecological concept by the work of C. H. Merriam with the U.S. Biological Survey of 1889. He described a series of life zones on the slopes of the San Francisco Peaks in Arizona, based on characteristic species of the flora and fauna. Other patterns of plant physiognomic and community zonation have now been cataloged in mountain ranges throughout the world. See LIFE ZONES. Merriam associated his life zones with temperature gradients present along mountain slopes. Later research on patterns of altitudinal zonation has centered on the response of species and groups of species to a complex of environmental gradients. Measurements of a species along a gradient, for example, the number of individuals, biomass, or ground coverage, generally form a bell-shaped curve. Peak response of a species occurs under optimum conditions and falls off at both ends of the gradient. The unique response of each species is determined by its physiological, reproductive, growth, and genetic characteristics. Zones of vegetation along mountain slopes are formed by intergrading combinations of species that differ in their tolerance to environmental conditions. Zones are usually indistinct entities rather than discrete groupings of species. However, under some conditions of localized disjunctions, very steep sections of gradients, or competitive exclusion, discon-
tinuities in the vegetation can create discrete communities. See ECOLOGICAL COMMUNITIES. Vegetation zones are often defined by the distributions of species having the dominant growth form, most frequently trees. For example, in the Cascade Mountains of Washington and Oregon, tundra above treeline has been called the sedge–grass zone, but successively lower zones are spruce–fir, arborvitae– hemlock, Douglas-fir, ponderosa pine, and finally several fescue–wheatgrass–sagebrush zones in which no trees grow. Another set of zones could be designated if other criteria were used. The boundaries of the middle elevation zones might be different if they had been defined instead by distributions of the dominant shrubs. Environmental gradients in mountains. Altitudinal vegetation zonation, therefore, is an expression of the response of individual species to environmental conditions. Plants along an altitudinal transect are exposed, not to a single environmental gradient, but to a complex of gradients, the most important of which are solar radiation, temperature, and precipitation. Although these major environmental gradients exist in most mountain ranges of the world, the gradients along a single altitudinal transect are not always smooth because of topographic and climatic variability. This environmental variation can result in irregular vegetation zones. The solar energy received by mountain surfaces increases with altitude, associated with decreases in air density and the amount of dust and water vapor. Global radiation levels in the European Alps have been shown to be 21% greater at 10,000 ft (3000 m) than at 650 ft (200 m) under a cloudless sky. An overcast sky is more efficient at reducing short-wave energy reaching low elevations and can increase the difference in energy input to 160%. However, more frequent clouds over high elevations relative to sunnier lower slopes commonly reduces this difference. See SOLAR RADIATION. Vegetation patterns are also strongly influenced by the decline in air temperature with increasing altitude, called the adiabatic lapse rate. Lapse rates are generally between 1.8◦F to 3.6◦F per 1000 ft (1◦C to 2◦C per 300 m), but vary with the amount of moisture present; wet air has a lower lapse rate. Thus, plants occurring at higher elevations generally experience cooler temperatures and shorter growing periods than low-elevation plants. Variation in the temperature gradient can be caused by differences in slope, aspect, radiation input, clouds, and air drainage patterns. See AIR TEMPERATURE. The precipitation gradient in most mountains is the reverse of the temperature gradient: precipitation increases with altitude. Moist air from low elevations is forced upward by the blocking action of the mountains into regions of cooler temperatures. Air holds less moisture at low temperatures than at warmer ones. When the atmosphere reaches the point of saturation (the dew point), condensation occurs, forming clouds or precipitation. This moisture gradient is most pronounced on windward slopes and in ranges close to warm and moist oceanic
Altitudinal vegetation zones
Europe
Southeast Asia
Tropical Andes
West Patagonia
East Himalaya Mexico
6000
altitude, m
Antarctica
New Zealand
5000
15,000
puna
4000 3000
sn
2000 tu
1000 0 North Pole
South Brazil Southeast Africa Southeast Australia
80°
70°
r nd
ow
e lin
a
60°
al
pi
ne
ve
ge
ta
tio
fe ni
rf
n
or
paramo t es
c l o u d f o re s t
d an
d f s an n ou s ee i d u re s t u o u rg r c de r fo ecid eve d d ad xe ife o mi con br
r bo
ea
lc
o
50°
40°
30°
e or
st
20°
10,000
t ro p i c a l mo n t a n e forest ical lowlan o r d t p rainforest rainforest
10°
0°
sub rai tropi nfo cal res t
10°
20°
sub
30°
ant
a rc
tic
40°
tus
ssnnoo ww llii ne soc kb elt
50°
60°
cool-temperate rainforest Fig. 1. Generalized patterns of altitudinal vegetation zonation in the Northern and Southern hemispheres. (After L. W. Price, Mountains and Man, University of California Press, 1981)
winds. As the moving air passes over the peaks of the mountain range and moisture in the air is depleted, precipitation is reduced, creating a rain “shadow” on the lee side of the range. See PRECIPITATION (METEOROLOGY). Temperate, tropical, and high latitudes. General changes in vegetation with increases in altitude include reduction in plant size, slower growth rates, lower production, communities composed of fewer species, and less interspecific competition. However, many regional exceptions to these trends exist. In western North America, for example, the lowest zone is often a treeless shrubland or prairie, so plant size initially increases with altitude as trees become important. Above this zone, the trend toward smaller size prevails. Characteristics of vegetation zones also vary with latitude (Fig. 1). Mountains at higher latitudes have predominantly seasonal climates, with major temperature and radiation extremes between summer and winter. Equatorial and tropical mountains have a strong diurnal pattern of temperature and radiation input with little seasonal variation. The upper altitudinal limit of trees, and the maximum elevation of plant growth generally, decreases with distance from the Equator, with the exception of a depression near the Equator. The following examples of altitudinal zonation in temperate, tropical, and high-latitude regions illustrate patterns that exist in each region. These patterns are not comprehensive, and many others exist locally and worldwide. Temperate. The vegetation of the Santa Catalina Mountains in Arizona illustrates the variety of zones found in a temperate mountain range (Fig. 2). Typically, precipitation increases with altitude, while air and soil temperatures decrease. Community diversity, the number of species present, decreases with altitude in all zones except the scrub desert, where arid conditions support fewer species. Biomass production is less at low elevations than on higher
slopes, the reverse of the general trend, because of limitations to growth under arid conditions. Hot and dry conditions at lower elevations support a treeless desert shrub and grassland, dominated by such shrubs as palo verde (Cercidium microphyllum) and mesquite (Prosopis juliflora). With increasing elevation, more moisture allows growth of broadleaf deciduous and evergreen oak trees (Quercus spp.). Needle-leaf evergreen trees, including ponderosa pine (Pinus ponderosa) and border limber pine (P. strobiformis), become increasingly important in the mid-elevation vegetation zones where temperatures are cooler, moisture more readily available, and the growing season longer. Broadleaf trees decrease in importance in these zones. Finally, in the upper zones, where the coolest temperatures and highest moisture levels prevail, conifers that commonly extend farther north including Douglasfir (Pseudostuga menziesii) and Engelmann spruce (Picea engelmanni) replace species better adapted to lower-elevation environments. At higher latitudes and where mountains are high enough, an alpine zone of grasses, forbs, cushion plants, and dwarf shrubs occurs above treeline. A similar sequential replacement of species in the altitudinal zones is repeated in both the herb and shrub layer. The Southern Hemisphere has fewer mountains in the temperate regions than areas north of the Equator. The oceanic climate of these mountains is cooler and moister because of the small land mass relative to the nearby oceans. A prolonged cold season is absent. In New Zealand, broadleaf evergreen trees such as evergreen beech (Nothofagus spp.) dominate moist mountain forests, with needle-leaf evergreens such as Podocarpus spp. being of lesser importance. Zones at all elevations have a richer, more diverse flora than their Northern Hemisphere counterparts. Within each zone, shrubs are more luxurious, and epiphytes and hanging vines are common. Tropical. Vegetation zones at all elevations in tropical mountains are dominated by a diurnal climate
5000 70°
80° South pole
altitude, ft
Spitsberg
529
Altitudinal vegetation zones
3000
subalpine forest
10,000
Community diversity, no. species
Temperature, °F (°C)
Precipitation, in. (cm)
9000
montane fir forest 2500 pine forest
8000
pine-oak forest
7000
2000 pine-oak woodland
1500
canyon woodland
open oak woodland
desert grassland 1000
sonoran desert scrub mesic
elevation, ft
elevation, m
530
48 (9)
80 (200)
13
55 (13)
65 (165)
18
63 (17)
50 (120)
58
72 (22)
40 (100)
33
6000
5000
4000
xeric
Fig. 2. Vegetation zones and gradients of mean annual soil temperature, annual precipitation, and community diversity on the south slopes of the Santa Catalina Mountains, Arizona, a temperate mountain range. (After R. H. Whittaker, and W. A. Niering, Vegetation of the Santa Catalina Mountains: A gradient analysis of the south slope, Ecology, 46: 429–451, 1965)
with little seasonality, because sun angles are high throughout the year. Growing seasons are determined by moisture availability rather than by temperature. At lower elevations in the Andes, tropical rainforests with a species-rich multilayered and luxurious canopy of broadleaf evergreen trees cover mountain slopes. These trees are accompanied by a lush mixture of hanging vines, epiphytes, shrubs, and flowering herbs. In the submontane and montane zones at higher elevations, the complexity of the forest is reduced. Trees are shorter and are adapted to drier conditions, and the canopy has fewer layers. Vines, epiphytes, and shrubs are present in fewer numbers, but mosses and lichens become more common. The subalpine zone is an elfin woodland or cloud forest of short, stunted trees covered with an abundance of mosses and lichens. Vegetation is adapted to the higher moisture and lower light levels created by clouds, which almost continuously envelop the zone. Plant communities in this zone are less complex than those lower on the mountain; fewer species are present and canopy structure is simpler. In the northern Andes of Colombia, the alpine zone consists of a low vegetation called paramo, which is dominated by arborescent members of the sunflower family, Compositae. Growth forms consist of tall tussock grasses, dwarf shrubs, herbs, and two forms unique to the tropics of the Southern Hemisphere: dendroid or tufted-leaf stemmed plants, and woolly candlelike plants. See PARAMO. Farther from the Equator in the central and southern Andes, the climate is drier and lower zones are
dominated by forests of deciduous trees or grassland savannas. The alpine zone consists of puna, a vegetation of smaller tussock grasses and many cushion and rosette plants. See PUNA; SAVANNA. High latitude. At high latitudes in both Northern and Southern hemispheres, fewer vegetation zones exist. Short growing seasons permit no tree growth, and shrubs and herbs are limited at higher elevations by a permanent cover of ice or snow, called the nival zone. Plant communities consist of fewer species than in temperate or tropical regions. In the Olgivie Mountains, Yukon Territory, Canada, valley bottoms are covered by willow (Salix spp.) and birch (Betula spp.) shrubs with a ground layer of grasses and forbs in drier sites, and tussock tundra dominated by cottongrass (Eriophorum spp.) in wet sites. At higher elevations, tall woody shrubs are absent, and this zone consists of a tundra of dwarf heath shrubs including mountain heather (Cassiope spp.) and mountain avens (Dryas spp.), grasses, forbs, and plants with cushion or rosette growth forms. See TERRESTRIAL ECOSYSTEM. John S. Campbell Bibliography. E. W. Beals, Vegetational change along altitudinal gradients, Science, 165:981-985, 1969; W. Lauer, The altitudinal belts of the vegetation in the central Mexican highlands and their climatic condition, Arctic Alpine Res., 5(pt. 2):A99-114, 1973; L. W. Price, Mountains and Man, 1981; C. Troll (ed.), Geoecology of the Mountainous Regions of the Tropical Americas, 1968; H. Walter, Vegetation of the Earth and Ecological Systems of the GeoBiosphere, 1979; R. H. Whittaker and W. A. Niering,
Aluminum Vegetation of the Santa Catalina Mountains: A gradient analysis of the south slope, Ecology, 46:429-451, 1965.
Alum A colorless to white crystalline substance which occurs naturally as the mineral kalunite and is a constituent of the mineral alunite. Alum is produced as aluminum sulfate by treating bauxite with sulfuric acid to yield alum cake or by treating the bauxite with caustic soda to yield papermaker’s alum. Other industrial alums are potash alum, ammonium alum, sodium alum, and chrom alum (potassium chromium sulfate). Major uses of alum are as an astringent, styptic, and emetic. For water purification alum is dissolved; it then crystallizes out into positively charged crystals that attract negatively charged organic impurities to form an aggregate sufficiently heavy to settle out. Alum is also used in sizing paper, dyeing fabrics, and tanning leather. With sodium bicarbonate it is used in baking powder and in some fire extinguishers. See ALUMINUM; COLLOID. Frank H. Rockett
Aluminum A metallic chemical element, symbol Al, atomic number 13, atomic weight 26.98154, in group 13 of the periodic system. Pure aluminum is soft and lacks strength, but it can be alloyed with other elements to increase strength and impart a number of useful properties. Alloys of aluminum are light, strong, and readily formable by many metalworking processes; they can be easily joined, cast, or machined, and accept a wide variety of finishes. Because of its many desirable physical, chemical, and metallurgical properties, aluminum has become the most widely used nonferrous metal. See PERIODIC TABLE. 1 1 H 3 Li 11 Na 19 K 37 Rb 55 Cs 87 Fr
2 4 Be 12 Mg 20 Ca 38 Sr 56 Ba 88 Ra
3 21 Sc 39 Y 71 Lu 103 Lr
4 22 Ti 40 Zr 72 Hf 104 Rf
lanthanide series actinide series
5 23 V 41 Nb 73 Ta 105 Db
6 24 Cr 42 Mo 74 W 106 Sg
7 25 Mn 43 Tc 75 Re 107 Bh
8 26 Fe 44 Ru 76 Os 108 Hs
9 27 Co 45 Rh 77 Ir 109 Mt
10 28 Ni 46 Pd 78 Pt 110 Ds
11 29 Cu 47 Ag 79 Au 111 Rg
12 30 Zn 48 Cd 80 Hg 112
13 5 B 13 Al 31 Ga 49 In 81 Tl 113
14 6 C 14 Si 32 Ge 50 Sn 82 Pb
15 16 7 8 N O 15 16 P S 33 34 As Se 51 52 Sb Te 83 84 Bi Po
57 58 59 60 61 62 63 64 65 La Ce Pr Nd Pm Sm Eu Gd Tb
66 67 Dy Ho
89 Ac
98 Cf
90 Th
91 Pa
92 93 94 95 96 97 U Np Pu Am Cm Bk
18 2 17 He 9 10 F Ne 17 18 Cl Ar 35 36 Br Kr 53 54 I Xe 85 86 At Rn
68 69 70 Er Tm Yb
99 100 101 102 Es Fm Md No
Aluminum is the most abundant metallic element on the Earth and Moon but is never found free in nature. The element is widely distributed in plants, and nearly all rocks, particularly igneous rocks, contain aluminum in the form of aluminum silicate minerals. When these minerals go intosolution, depend-
Physiochemical properties of pure aluminum Property
Value
Atomic number Atomic weight Crystal structure Density, g/cm3 at 25°C Melting point Boiling point Latent heat of fusion, cal/g Latent heat of vaporization, cal/g Heat of combustion, cal/g at 25°C Specific heat, cal/g/°C at 25°C Thermal conductivity, cal/cm/°C/s at 25°C Coefficient of thermal expansion, 10⫺6 /°C Thermal diffusivity, cm2 /s at 25°C Electrical resistivity, microhm-cm at 25°C Hardness, Mohs Reflectivity
13 26.98154 Face-centered cubic 2.698 660.37°C (1220.7°F) 2447°C (4436.6°F) 94.9 2576 7420 0.215 0.566 23 0.969 2.7 2–2.9 85–90%
ing upon the chemical conditions, aluminum can be precipitated out of the solution as clay minerals or aluminum hydroxides, or both. Under such conditions bauxites are formed. Bauxites serve as principal raw materials for aluminum production. See BAUXITE; CLAY MINERALS; IGNEOUS ROCKS; WEATHERING PROCESSES. Aluminum is a silvery metal. Naturally occurring 27 aluminum consists of a single isotope, 13 Al. Aluminum crystallizes in the face-centered cubic structure with edge of the unit lattice cube of 4.0495 angstroms. Aluminum is known for its high electrical and thermal conductivities and its high reflectivity. A summary of some of the important properties of pure aluminum is given the table. The electronic configuration of the element is 1s22s22p63s23pl. Aluminum exhibits a valence of +3 in all compounds, with the exception of a few high-temperature monovalent and divalent gaseous species. Aluminum is stable in air and resistant to corrosion by seawater and many aqueous solutions and other chemical agents. This is due to protection of the metal by a tough, impervious film of oxide. At a purity greater than 99.95%, aluminum resists attack by most acids but dissolves in aqua regia. Its oxide film dissolves in alkaline solutions, and corrosion is rapid. See CORROSION. Aluminum is amphoteric and can react with mineral acids to form soluble salts and to evolve hydrogen. Molten aluminum can react explosively with water. The molten metal should not be allowed to contact damp tools or containers. At high temperatures aluminum reduces many compounds containing oxygen, particularly metal oxides. These reactions are used in the manufacture of certain metals and alloys. Applications in building and construction represent the largest single market of the aluminum industry. Millions of homes use aluminum doors, siding,
531
532
Aluminum (metallurgy) windows, screening, and down-spouts and gutters. Aluminum is also a major industrial building product. Transportation is the second largest market. In automobiles, aluminum is apparent in interior and exterior trim, grilles, wheels, air conditioners, automatic transmissions, and some radiators, engine blocks, and body panels. Aluminum is also found in rapid-transit car bodies, rail cars, forged truck wheels, cargo containers, and in highway signs, divider rails, and lighting standards. In aerospace, aluminum is found in aircraft engines, frames, skins, landing gear, and interiors. The food packaging industry is a fast-growing market. In electrical applications, aluminum wire and cable are major products. Aluminum appears in the home as cooking utensils, cooking foil, hardware, tools, portable appliances, air conditioners, freezers, and refrigerators. There are hundreds of chemical uses of aluminum and aluminum compounds. Aluminum powder is used in paints, rocket fuels, and explosives, and as a chemical reductant. See ALUMINUM ALLOYS. Allen S. Russell Bibliography. F. A. Cotton, et al., Advanced Inorganic Chemistry, 6th ed., Wiley-Interscience, 1999; D. R. Lide, CRC Handbook Chemistry and Physics, 85th ed., CRC Press, 2004.
Aluminum (metallurgy) The separation, extraction, and purification of alumina from ores, followed by the production of aluminum. Aluminum is the most abundant metallic element on the Earth and Moon but is never found free in nature. It makes up more than 8% of the solid portion of the Earth’s surface. Seawater contains an average of only 0.5 ppm aluminum. The element is widely distributed in plants, where it may be present in significant concentrations, particularly in vegetation in marshy places and acid soils. Nearly all rocks, particularly igneous rocks, contain aluminum in the form of aluminum silicate minerals. When subjected to weathering under the right conditions, these minerals go into solution; then, depending upon the chemical conditions, aluminum can be precipitated out of the solution as clay minerals or aluminum hydroxides, or both. Any quartz in the parent rock is relatively unattacked by the weathering processes and remains as sand in the deposit. Under conditions which precipitate clay minerals and leach out iron and the alkali and alkalineearth metals, deposits of high-grade kaolin clay can result. Sometimes natural elutriation or mechanical separation by washing occurs, which separates the clay from the sand and enhances the purity of the deposit. Under conditions which precipitate aluminum principally as hydroxide minerals and leach out alkali and alkaline-earth metals completely and iron to varying degrees, bauxites are formed. It is also possible for clay minerals to be desilicified to aluminum hydroxides and for aluminum hydroxides to be resilicified to clay minerals. See BAUXITE; CLAY MINERALS; IGNEOUS ROCKS; WEATHERING PROCESSES.
Bauxites usually consist of mixtures of the following minerals in varying proportions: gibbsite, also known as hydrargillite [Al(OH)3]; boehmite and diaspore [AlO(OH)]; clay minerals such as kaolinite [Al2Si2O5(OH)4]; quartz [SiO2], and anatase, rutile, and brookite [TiO2]. Small amounts of magnetite [Fe3O4], ilmenite [FeTiO3], and corundum [Al2O3] are sometimes present. In addition to the above minerals, bauxites usually contain traces of other insoluble oxides such as zirconium, vanadium, gallium, chromium, and manganese. The term alumina trihydrate is often applied to the mineral gibbsite, and the term alumina monohydrate to the minerals boehmite and diaspore. However, the minerals are not hydrates in the true sense of the word. However, for processing purposes, bauxites are often classified as trihydrate bauxites and monohydrate bauxites, depending upon their content of the principal mineral or minerals containing the extractable alumina. In general, European bauxites are of the monohydrate type and are geologically older than bauxites in tropical countries, which occur generally as the trihydrate type. This, however, is not an absolute rule. Some European bauxites contain gibbsite along with boehmite and diaspore, and Caribbean bauxites contain small-to-medium amounts of boehmite along with the gibbsite. Typical bauxites which are used for aluminum production contain 40–60% total alumina (Al2O3), 1–15% total silica (SiO2), 7–30% total Fe2O3, 1–5% titania (TiO2), and 12–30% combined H2O. The suitability of a bauxite as a raw material for aluminum production depends not only on its alumina content but also on its content of combined silica, which is usually in the form of the mineral kaolinite. Kaolinite not only contains aluminum that cannot be extracted in the Bayer process, but also reacts with the sodium–aluminate solution to cause a loss of sodium hydroxide. In some bauxites, the clay minerals are concentrated in the fine-particle size range. In this case, if the economics are favorable, the bauxite can be beneficiated by washing it on screens to remove substantial amounts of the clay and thus make a product containing higher alumina and lower combined silica contents than the original bauxite (Fig. 1). History. H. C. Oersted probably prepared the first metallic aluminum (impure) in 1824 by reducing aluminum chloride with potassium amalgam. F. W¨ ohler is generally credited with the discovery of the metal by virtue of the first isolation of less impure aluminum in 1827 and the first description of a number of its properties. The metal remained a laboratory curiosity until 1854, when Henri Sainte-Claire Deville improved the earlier method of preparation by employing sodium as reductant and was successful in producing larger quantities of relatively pure metal. The first electrolytic preparation of aluminum, which also occurred in 1854, was accomplished by electrolysis of fused sodium aluminum chloride independently in the laboratories of Deville in France and of R. Bunsen in Germany. See ELECTROLYSIS. The modern industrial electrolytic method of production was discovered simultaneously and independently by Charles Martin Hall in the United States
Aluminum (metallurgy)
Fig. 1. Bauxite washing. (Aluminum Company of America)
and Paul-Louis Héroult of France in 1886. The essentials of their discovery remain the basis for today’s aluminum industry. Alumina dissolved in a fluoride fusion (largely cryolite, Na3AlF6) is electrolyzed with direct current. Carbon dioxide is discharged at the anode, while the metal is deposited on molten aluminum lying on the carbon lining at the bottom of the cell. The technology for extracting aluminum from its ores was further improved in 1888, when Karl Josef Bayer patented in Germany a method for making pure alumina from bauxite. Production. The most widely used technology for producing aluminum involves two steps: extraction and purification of alumina from ores, and electrolysis of the oxide after it has been dissolved in fused cryolite. Bauxite is by far the most used raw material. Technically feasible processes exist for extracting alumina from other raw materials such as clays, anorthosite, nepheline syenite, and alunite; however, these processes have not been competitive with the Bayer process in the United States. With the present rapidly changing economics of bauxite production, some of these processes may become competitive in the future. A few of these alternate raw materials are used as a source for alumina in Europe and Asia. Alumina extraction. A general flow sheet for the Bayer process is shown in Fig. 2. In the process, bauxite is crushed and ground, then digested at elevated temperature (280–450◦F or 140–230◦C) and pressure in a strong solution of caustic soda (80–110 g Na2O/liter
or 0.7–0.9 lb Na2O/gal). For monohydrate-type bauxites, in which the alumina occurs in forms which are more difficult to dissolve than in trihydrate-type bauxites, stronger solutions (up to 220 g Na2O/liter or 1.8 lb Na2O/gal), higher temperatures (up to 570◦F or 300◦C) and pressures (as high as 150 atm or 1.52 × 107 pascals), and sometimes longer digestion times are required. The gibbsite, boehmite, or diaspore in the bauxites reacts with the sodium hydroxide to form soluble sodium aluminate. The residue, known as red mud, contains the insoluble impurities and the sodium aluminum silicate compound, referred to as desilication product, formed by the reaction of clay minerals with the sodium aluminate–sodium hydroxide solution. The red mud is separated from the solution by countercurrent decantation and filtration. After cooling, the solution is supersaturated with respect to alumina. It is seeded with recycled synthetic gibbsite (alumina trihydrate) and agitated. A large part of the alumina in solution thus crystallizes out as gibbsite. This gibbsite is classified into product and seed, the seed being recycled and the product washed. Wash water and spent liquor, after concentration by evaporation, are recycled to the digestion system. For metal production, the product is calcined at temperatures up to 2400◦F (1300◦C) to produce alumina containing about 0.3–0.8% soda, less than 0.1% iron oxide plus silica, and trace amounts of other oxides. Some of the precipitated gibbsite is sold for use in the chemicals industry. Specially precipitated pigment-grade gibbsite is used for rubber,
533
Aluminum (metallurgy) storage
bauxite from mine
rod mill
steam crushed lime
waste lake
filters heat exchangers
digesters green liquor
crushed bauxite
washing thickeners
thickeners
sand classifier
to smelting
precipitators
blow-off tanks
primary classifying thickener
slurry mixer
alumina (white powder)
heaters caustic
d
heaters
spent liquor flow
se e
534
cooler rotary kiln
tertiary classifiers
secondary classifying thickener
filter
Fig. 2. Processing of alumina from bauxite. (Aluminum Company of America)
plastic and paper fillers, and paper coating. Activated aluminas having high internal surface area are made from the gibbsite by calcination at low-to-moderate temperatures. They are used as desiccants and catalysts. Fully calcined aluminas are used in the production of ceramics, abrasives, and refractories. Alumina can be recovered from low-grade bauxites containing high concentrations of combined silica by the so-called combination or lime–soda sinter process. The process consists of treatment by the Bayer process, following which the red mud is mixed with calcium oxide and sodium carbonate and calcined. This treatment forms sodium aluminate from the aluminous phases in the red mud. The sodium aluminate is then leached out with water, and alumina is recovered from the resulting solution. Processes for treating other ores have been developed which consist of calcination or fusion of the ore with limestone to make calcium aluminate, from which alumina is leached out with sodium carbonate solution and recovered. Such a process is used in Russia for treating nepheline syenite to yield alumina, and cement as a by-product. High-iron bauxites have been smelted with limestone in the Pedersen process to yield pig iron and a calcium aluminate slag that can be leached with sodium carbonate solution to recover alumina. Electrolytic reduction (smelting). Although unchanged in principle, the smelting process of today differs in detail and in scale from the original process discovered by Hall and Héroult. Technology has effected substantial improvements in equipment, materials, and control of the process, and has lowered the energy and labor requirements and the final cost of primary metal. In a modern smelter (Fig. 3), alumina is dissolved in cells (pots)—rectangular steel shells lined with
carbon—containing a molten electrolyte (bath) consisting mostly of cryolite. The bath usually contains 2–8% alumina. Excess aluminum fluoride and calcium fluoride are added to lower the melting point and to improve operation. Carbon anodes are hung from above the cells with their lower ends extending to within about 1.5 in. (3.8 cm) of the molten metal, which forms a layer under the molten bath. The heat required to keep the bath molten is supplied by the electrical resistance of the bath as current passes through it. The amount of heat developed with a given current depends on the length of the current path through the electrolyte, that is, anode– cathode distance, which is adjusted to maintain the desired operating temperature, usually 1760–1780◦F (960–970◦C). A crust of frozen bath, 1–3 in. (2.5– 7.5 cm) thick, forms on the top surface of the bath and on the sidewalls of the cell. Alumina is added to the bath or on the crust, where its sorbed moisture is driven off by heat from the cell. While preheating on the crust, the alumina charge serves as thermal insulation. Periodically, the crust is broken and the alumina is stirred into the bath to maintain proper concentration. See CRYOLITE. The passage of direct current through the electrolyte decomposes the dissolved alumina. Metal is deposited on the cathode, and oxygen on the gradually consumed anode. About 0.5 lb (0.23 kg) of carbon is consumed for every pound of aluminum produced. The smelting process is continuous. Alumina is added, anodes are replaced, and molten aluminum is periodically siphoned off without interrupting current to the cells. Current efficiencies in the industrial electrolytic process are about 85–92%, and the energy efficiency is about 40%. The voltage at the cell terminals is 4–6 V, depending on the size and condition of the
Aluminum (metallurgy) cell. Voltage is required to force the current through the entire cell, and the corresponding power (voltage × amperage) is largely converted into heat in the bath. The amount of power required to maintain the temperature is a smaller proportion of the total power input in large cells than in small ones because of the lower ratio of surface to volume. Thus, power consumed per pound of metal is somewhat less in large than in small cells. Consumption of 6–8 kWh/lb (13–18 kWh/kg) of aluminum produced includes bus-bar, transformer, and rectifier losses. A potline may consist of 50–200 cells with a total line voltage of up to 1000 V at current loads of 50,000–225,000 A. Electric power is one of the most costly raw materials in aluminum production. Aluminum producers have continually searched for sources of cheap hydroelectric power, but have also had to construct facilities that produce power from fossil fuels. In the past half century, technological advances have significantly reduced the amount of electrical energy necessary to produce a pound of aluminum. Current is led out of the cell to the anode bus-bar by a number of carbon block anodes suspended in parallel rows on vertical conducting rods of copper or aluminum. Because impurities in the anodes dissolve in the bath as they are consumed, pure carbon
(calcined petroleum coke or pitch coke) is used as raw material. The ground coke is mixed hot with enough coal tar or petroleum pitch to bond it into a block when pressed in a mold to form the “green anode.” This is then baked slowly at temperatures up to 2000–2200◦F (1100–1200◦C). In a cavity molded in the top of each block, a steel stub is embedded by casting molten iron around it or by using a carbonaceous paste; the conducting bar is bolted to this stub. Such an electrode is termed a prebaked anode to distinguish it from the Soderberg anode, in which the electrode (single large anode to a cell) is formed in place from a carbonaceous paste that is baked by heat from the pot as it gradually descends into the electrolyte. The steel shell of the pot is thermally insulated and lined with carbon. The carbon bottom, covered with molten aluminum, serves as cathode of the cell. Electrical connection is made to the carbon cathode by steel bars running through the base of the cell and embedded in the carbon lining. Molten aluminum is siphoned from the smelting cells into large crucibles. From there the metal may be poured directly into molds to produce foundry ingot, or transferred to holding furnaces for further refining or alloying with other metals to form fabricating ingot. As it comes from the cell, primary aluminum averages about 99.8% purity.
direct current
alumina supply hopper anode adjusting wheel
carbon anodes
signal light
alumina hopper
measuring gate
electrical insulation reduction pots in series current
bus bar stirring
steel shell heat insulation carbon lining (cathode)
tapping holding furnace
molten aluminum
cryolite bath reduction pot
535
charging to fabrication
aluminum ingot Fig. 3. Production of aluminum metal. (Aluminum Company of America)
casting
536
Aluminum (metallurgy) See ELECTROMETALLURGY; PYROMETALLURGY, NONFERROUS. Melting. In plants not adjacent to a smelter, it is necessary to remelt charges, usually consisting of mill scrap or returned scrap with enough primary metal to provide composition control. For casting ingot for fabrication, gas- or oil-fired reverberatory furnaces are commonly used. Molten metal is generally transferred from the melting furnace to a reverberatory holding furnace or to a ladle for subsequent casting. Metal for foundry use in producing sand, permanent mold, and die castings is usually remelted in gas- or oil-fired crucible furnaces, although in large-scale operations reverberatory furnaces may be used. Scrap recycling has grown through the years to considerable economic importance, and provides a valuable source of aluminum at a much lower energy expenditure than for primary metal. For many years, secondary producers have purchased scrap and reclaimed it, generally into foundry ingot for remelting. Primary producers have also purchased large amounts of mill scrap, as well as discarded cans, and recycled it into ingot used to fabricate sheet for more cans. This closed-circuit type of recycling saves energy, helps to preserve the environment, and conserves metal to the greatest degree. After remelting, the molten aluminum alloy must be treated to remove dissolved hydrogen, inclusions such as oxides formed during remelting or from the surface of the charge material, and undesirable trace elements, such as sodium, in some alloys. Classically, the practice has been to bubble chlorine or a mixture of nitrogen and chlorine through the melt by means of graphite tubes. Modern practice makes use of a wide variety of in-line systems in which the metal is treated by filtering or fluxing, or both, during transfer between the holding furnace and the casting unit. These processes are more efficient than furnace fluxing, and their use minimizes air pollution, improves metal quality, and saves production time. Although many ingot casting methods have been used historically, today most ingots for fabricating are cast by the direct chill process or some modification thereof. This process is accomplished on a semicontinuous basis in vertical casting and on a continuous basis in horizontal casting. The direct chill process consists of pouring the metal into a short mold. The base of the mold is a separate platform that is lowered gradually into a pit while the metal is solidifying. The frozen outer shell of metal retains the stillliquid portion of the metal when the shell is past the mold wall. Water is used to cool the mold, and water impingement on the ingot shell also cools the ingot as it is lowered. Ingot length is limited only by the depth to which the platform can be lowered. In horizontal casting, the same principles are applied, but the mold and base are turned so that the base moves horizontally, removing the constraint of pit depth. Much longer ingot can be practically cast by the horizontal method. By sawing the ingot while casting is progressing, the process can be made continuous.
Fabrication methods. Aluminum alloys are fabricated by all the methods known to metalworking and are commercially available in the form of castings; in wrought forms produced by rolling, extruding, forging, or drawing; in compacted and sintered shapes made from powder; and in other monolithic forms made by a wide variety of processes. Some of the common fabrication processes are represented in Fig. 4. In addition, pieces produced separately can be assembled by common joining procedures to build up complex shapes and structures. See ALUMINUM ALLOYS; METAL FORMING. Casting. In casting aluminum alloys, molten metal is poured, or forced by pressure, into a mold where it solidifies. The temperature is usually 1150–1400◦F (620–760◦C) depending on the alloy composition and on the type of casting. For die casting and permanent mold casting, the molds are metal and are used repeatedly. Other processes use expendable molds made of sand or plaster. See METAL CASTING. Rolling. Rolling is the process of reducing material by passing it between pairs of rolls. Aluminum alloy rolled products include plate (0.250 in. or 0.6 cm or more thick), sheet (0.006–0.249 in. or 0.02– 0.632 cm thick), and foil (less than 0.006 in. thick). The starting product is ingot, which may be up to 360 in. (914 cm) long, 72 in. (182 cm) wide, and 26 in. (66 cm) thick. Initial reduction is hot (800– 1100◦F, or 426–593◦C, depending on alloy) and is done in a reversing mill. Subsequent rolling operations may be in multistand continuous mills or singlestand mills, and may be hot or cold depending on the thickness and properties required in the final product. Other aluminum alloy rolled products include rod and bar, which are passed through grooved rolls. See METAL ROLLING. Extruding. In extruding aluminum alloys, an ingot or billet is forced by high pressure to flow from a container through a die opening to form an elongated shape or tube, usually at metal temperatures between 550 and 1050◦F (287 and 565◦C). Hydraulic presses with capacities of 500 to 14,000 tons (4.5 × 106 to 1.2 × 108 newtons) are used. The extrusion process is capable of producing sections with weights of a few ounces to more than 200 lb/ft (300 kg/m), with thicknesses from a few tenths of an inch to about 10 in. (25 cm), with circumscribing circle diameters of 0.25 in. (0.64 cm) to about 3 ft (0.9 m), and lengths in excess of 100 ft (30 m). Tubing in diameters of 0.25–33 in. (0.64–0.84 cm) and pipe up to 20 in. (50 cm) are also produced. Stock that is produced by extrusion is also fabricated to tubing by drawing and to forgings by subsequent metalworking operations. See EXTRUSION. Forging. Forgings are produced by pressing or hammering aluminum alloy ingots or billets into simple rectangular or round shapes on flat dies and into complex forms in cavity dies. Hydraulic presses capable of forces up to 75,000 tons (6.60 × 108 N) and mechanical presses with capacities up to 16,000 tons are used. Forging hammers have ram weights of 500–110,000 lb (227–50,000 kg). Metal temperatures of 600–880◦F (315–471◦C) are generally used,
Aluminum (metallurgy)
remelting furnace
dc ingot caster
hot reversing mill
end shear
heating furnace
heat-treating coldfurnace rolling mill shear
edge shear
scalper
537
coiler continuous mill
quench tank
annealing furnace
aging furnace
finish shear
sheet and plate
stretcher
reducing mill
foil foil spooling mill mill foil
coil annealing furnace
remelting furnace
heating furnace
dc ingot caster
heatfinish treating quench straightener rolling furnace chamber mill draw bench saw
scalper
blooming mill
aging furnace bar and rod
wirecoiler annealing drawing furnace finish rolling mill machine
reheat furnace
wire
alloying materials aluminum ingot aluminum scrap
extrusion press
remelting saw furnace
heat-treating furnace
aging furnace saw
stretcher heating furnace
dc ingot caster remelting furnace
heating furnace
annealing furnace saw
draw bench
heat-treating furnace saw stretcher
pointer extrusion press
Fig. 4. Fabrication operations. (Aluminum Company of America)
extruded shapes and tube
quench chamber
saw
dc hollow ingot caster
contour roller
quench tank
aging furnace
drawn tube
538
Aluminum alloys depending on alloy and type of forging. See FORGING. Drawing. Aluminum alloys can be made into deepdrawn shapes, of which pots and pans and beverage cans are common examples, by forcing sheet or foil into deep holes or cavities of the desired shape. Other drawn products, such as wire and tubing, are produced by pulling rolled or extruded stock through dies. See DRAWING OF METAL. Compacting and sintering. Aluminum powders produced by atomizing can be compacted, either alone or mixed with powders of other elements, in dies of complex shape. The compact is then heated to a high temperature to promote additional bonding between the particles to result in products with desired shapes. See POWDER METALLURGY. Other processes. Aluminum products are machined to final shape and dimension on lathes, joiners, routers, drill presses, shears, grinders, and many other high-speed metal-cutting and metalworking machines. Other fabricating techniques include impact extrusion, forming, stamping, embossing, coining, bending, rotary swaging, and cold heading. Joining. Aluminum products formed into shapes by any of the above processes can be joined by welding, brazing, soldering, and adhesive bonding, and by mechanical means such as crimping, seaming, screwing, bolting, and riveting. Products and uses. Applications in building and construction represent the largest single market of the aluminum industry. Millions of homes use aluminum doors, siding, windows, screening, and downspouts and gutters, which require little maintenance and provide a long life. Aluminum is also a major industrial building product. Excellent weather resistance makes it ideal for all climates and locations. Selected aluminum alloys are suitable near the seacoast, where salt spray may be deleterious to other metals. Colored and given additional protection by an electrolytic process called anodizing, aluminum appears in the curtain wall construction of many of the world’s tallest buildings. Transportation is the second largest market. In automobiles, aluminum is apparent in interior and exterior trim, grilles, wheels, and air conditioners. Another major use is in automatic transmissions. Some auto radiators, engine blocks, and body panels are made of aluminum to hold down weight and improve fuel economy. Aluminum is also found in rapidtransit car bodies, engine parts for diesel locomotives, rail freight and tank cars, bus and truck engines, forged truck wheels, cargo containers, and in highway signs, divider rails, and lighting standards. Aluminum’s light weight and corrosion resistance are responsible for applications in pleasure-boat hulls, fishing boats, tanks for ships transporting liquefied natural gas, and deck houses for naval vessels. In aerospace, aluminum is found in aircraft engines, frames, skins, landing gear, and interiors. Toughness, lightness, and heat-reflective characteristics have made it the preferred material for satellites and other spacecraft. See AIRFRAME; SPACECRAFT STRUCTURE.
In the packaging industry, aluminum helps in the preparation of foods and beverages and keeps them pure during distribution and storage. The fastcooling, easy-opening, lightweight, and recyclable features of the all-aluminum can have resulted in the use of billions of aluminum beverage containers. Foil pouches and bags, twist-off closures, and easy-open ends revolutionized the food and beverage packaging industries. See FOOD ENGINEERING. In electrical applications, aluminum wire and cable are major products. Underground electrical cables require large amounts of aluminum. Aluminum wiring is also used in residential, commercial, and industrial buildings. See TRANSMISSION LINES; WIRING. Consumer use began before the 1900s with aluminum cooking utensils. In modern homes, aluminum appears as cooking foil, hardware, tools, portable appliances, air conditioners, freezers, and refrigerators. There are hundreds of chemical uses of aluminum and aluminum compounds. Aluminum powder is used in paints, rocket fuels, and explosives, and as a chemical reductant. Allen S. Russell Bibliography. Aluminum Corporation of America, Melting and Casting Aluminum, 1995; G. S. Brady, H. R. Clauser, and J. A. Vaccari, Materials Handbook, 15th ed., 2002; J. E. Hatch (ed.), Aluminum: Properties and Physical Metallurgy, 1984; G. E. Totten and D. S. MacKenzie (eds.), Handbook of Aluminum, vol. 1: Physical Metallurgy and Processes, 2003.
Aluminum alloys Substances formed by the addition of one or more elements, usually metals, to aluminum. The principal alloying elements in aluminum-base alloys are magnesium (Mg), silicon (Si), copper (Cu), zinc (Zn), and manganese (Mn). In wrought products, which constitute the greatest use of aluminum, the alloys are identified by four-digit numbers of the form NXXX, where the value of N denotes the alloy type and the principal alloying element(s) as follows: 1 (Al; at least 99% aluminum by weight), 2 (Cu), 3 (Mn), 4 (Si), 5 (Mg), 6 (Mg + Si), 7 (Zn), 8 (other). See COPPER; MAGNESIUM; MANGANESE; SILICON; ZINC. Iron and silicon are commonly present as impurities in aluminum alloys, although the amounts may be controlled to achieve specific mechanical or physical properties. Minor amounts of other elements, such as chromium (Cr), zirconium (Zr), vanadium (V), lead (Pb), and bismuth (Bi), are added to specific alloys for special purposes. Titanium additions are frequently employed to produce a refined cast structure. See IRON; TITANIUM. Aluminum-base alloys are generally prepared by making the alloying additions to molten aluminum, forming a liquid solution. As the alloy freezes, phase separation occurs to satisfy phase equilibria requirements and the decrease in solubility as the temperature is lowered. The resultant solidified structure consists of grains of aluminum-rich solid solution
Aluminum alloys TABLE 1. Nominal composition and forming processes for common wrought aluminum alloys Alloy 1100 1350 2011 2014 2024 2036 2219 3003 3004 5052 5083 5086 5182 5657 6061 6063 7005 7050 7075 *
Form*
Si
Cu
Mg
Mn
Zn
Other
b–d b–d b–d b–e b–e b b, e b–d b–d b–d b, e b–d b b, c b–e b–e b–e b, e b–e
— — — 0.8 — — — — — — — — — — 0.6 0.4 — — —
— — 5.5 4.4 4.4 2.6 6.3 — — — — — — — 0.25 — — 2.4 1.6
0.1 — — 0.4 1.5 0.45 — — 1.0 2.5 4.5 4.0 4.5 0.8 1.0 0.7 1.5 2.3 2.5
— — — 0.8 0.6 0.25 0.3 1.2 1.2 — 0.7 0.5 0.35 — — — 0.5 — —
— — — — — — — — — — — — — — — — 4.5 6.2 5.6
99.00 Al (min.) 99.5 Al (min.) 0.5 Pb, 0.5 Bi — — — 0.1 V, 0.1 Zr — — 0.25 Cr 0.15 Cr 0.15 Cr — — 0.25 Cr — 0.15 Cr, 0.14 Zr 0.12 Zr 0.25 Cr
b ⫽ rolled; c ⫽ drawn; d ⫽ extruded; and e ⫽ forged.
and crystals of intermetallic compounds. Elements which lower the freezing point of aluminum, such as copper, magnesium, and silicon, tend to segregate to the portions of the grains which freeze last, such as the cell boundaries. Elements which raise the freezing point, such as chromium or titanium (Ti), segregate in the opposite manner. See EUTECTICS; SOLID SOLUTION. A decrease in solubility with falling temperature also provides the basis for heat treatment of solid aluminum alloys. In this operation, the alloy is held for some time at a high temperature to promote dissolution of soluble phases and homogenization of the alloy by diffusion processes. The limiting temperature is the melting point of the lowest melting phase present. The time required depends both on temperature and on the distances over which diffusion must occur to achieve the desired degree of homogenization. Times of several hours can be necessary with coarse structures such as sand castings. Only a few minutes may be adequate, however, for rapidly heated thin sheet. See HEAT TREATMENT (METALLURGY). The solution heat treatment is followed by a quenching operation in which the article is rapidly cooled, for example, by plunging it into cold or hot water or by the use of an air blast. This produces a supersaturated metallic solid solution that is thermodynamically unstable at room temperature. In several important alloy classes, such as 2XXX, 6XXX, and 7XXX, the supersaturated solution decomposes at room temperature to form fine, submicroscopic segregates or precipitates that are precursors to the equilibrium phases predicted by phase diagrams. The precipitation phenomenon, occurring over periods of days to years, produces substantial increases in strength. Additional precipitation strengthening can be obtained by heating the alloy at temperatures in the range 250–450◦F (121–232◦C), the time and temperature varying with alloy composition and the objectives with respect to mechanical
properties and other characteristics such as corrosion resistance. Wrought alloys. Table 1 lists the nominal compositions of a number of commercially important wrought alloys and the type of products for which they are used. The alloys are generally classified in two broad categories depending upon their response to heat treatment as described in the preceding paragraph. Those having no or minor response are identified as non-heat-treatable alloys and include the 1XXX, 3XXX, and 5XXX compositions. Those that do respond are known as heat-treatable alloys and include the 2XXX, 6XXX, and 7XXX compositions. Included in the non-heat-treatable group is 1350, a special grade used for electrical conductor products. Alloy 1100 is a grade of 99.0% minimum aluminum content with particular controls on iron (Fe), silicon, and copper contents, available in a variety of product forms such as sheet, foil, wire, rod, and tube, used for packaging, fin stock, and a variety of sheet metal applications. The manganesecontaining alloy 3003 is a moderate-strength, very workable alloy for cooking utensils, tube, packaging, and lithographic sheet applications. The stronger aluminum-manganese-magnesium alloy 3004 is used for architectural applications, for storage tanks, and especially for drawn and ironed beer and beverage containers. Alloy 5052 is a workable, corrosion-resistant aluminum-magnesium alloy for many metalworking and metal-forming purposes and marine applications. Where higher strength is required, the higher-Mg-content, weldable 5086 or 5083 alloys may be used. The latter is employed in construction of welded tanks for liquefied gas (cryogenic) transport and storage and for armor plate in military vehicles. Alloy 5182 is also a high-strength aluminum-magnesium alloy that is employed primarily in a highly strain-hardened condition for beverage can ends. Alloy 5657 is a lower-strength material
539
Aluminum alloys TABLE 2. Nominal composition and casting procedure for common aluminum casting alloys Alloy
Form∗
Si
413.0 B† 443.0 † F 332.0 355.0 356.0 380.0 390.0
D B, C C B, C B, C D C, D
12 5.3 9.5 5.0 7.0 8.5 17
Cu
.15 max 3.0 1.3 3.5 4.5
Mg
1.0 0.5 0.3 0.55
∗
B ⫽ sand casting; C ⫽ permanent mold casting; and D ⫽ die casting. letter indicates modifications of alloys of the same general composition or differences in impurity limits, from alloys having the same four-digit numerical designations. † The
produced with a bright anodized finish for automobile trim and other decorative applications. The heat-treatable alloys 2014, 2024, and 7075 have high strengths and are employed in aircraft and other transportation applications. Modifications of these basic alloys, such as 2124 and 7475, were developed to provide increased fracture toughness. High toughness at high strength levels is achieved in thick-section products with 7050. Where elevated temperatures are involved, the 2XXX alloys are preferred. One such alloy, 2219, also has good toughness at cryogenic temperatures and is weldable. This alloy was prominently employed in the fuel and oxidizer tanks serving as the primary structure of the Saturn space vehicle boosters. Alloy 6061 is used for structural applications where somewhat lower strengths are acceptable. For example, 6061 may be used for trailers, trucks, and other transportation applications. Alloy 6063 is a stilllower-strength, extrudable, heat-treatable alloy for furniture and architectural applications. The use of lead and bismuth in 2011 produces a heat-treatable, free-machining alloy for screw-machine products. Alloy 2036 is a moderate-strength alloy with good workability and formability, employed as body sheet in automotive applications. Casting alloys. Table 2 shows the nominal compositions of several important casting alloys. The major alloying addition is silicon, which improves the castability of aluminum and provides moderate strength. Other elements are added primarily to increase the tensile strength. Most die castings are made of alloy 413.0 or 380.0. Alloy 443.0 has been very popular in architectural work, while 355.0 and 356.0 are the principal alloys for sand castings. Number 390.0 is employed for die-cast automotive engine cylinder blocks, while alloy F332.0 is used for pistons for internal combustion engines. Casting alloys are significant users of secondary metal (recovered from scrap for reuse). Thus, casting alloys usually contain minor amounts of a variety of elements; these do no harm as long as they are kept within certain limits. The use of secondary metal is also of increasing importance in wrought alloy manufacturing as producers take steps to reduce the energy that is required in producing fabricated aluminum products. See ALLOY; ALLOY STRUCTURES; METAL CASTING. Allen S. Russell
High-strength alloys. Since aluminum comprises 70–80% of the weight of an airframe, metallurgists have been pursuing aluminum alloy development programs directed toward producing materials which would be characterized by stronger, stiffer, and lighter-weight properties. In addition, the titanium alloys of aircraft gas turbines are prime targets for replacement by lighter-weight alloys. Researchers have improved aluminum alloys by using techniques that add lithium and by using three processing technologies comprising powder metallurgy, mechanical alloying, and a process which involves blending aluminum alloy powders and silicon carbide fibers to form a composite. Aluminum-lithium alloys. Lithium is the lightest metal in existence. For each weight percent of lithium added to aluminum, there is a corresponding decrease of 3% in the alloy’s weight. Therefore, aluminum-lithium alloys offer the attractive property of low density. A second beneficial effect of lithium additions is the increase in elastic modulus (stiffness). Also, as the amount of lithium is increased, there is a corresponding increase in strength due to the presence of very small precipitates which act as strengthening agents to the aluminum. As the precipitates grow during heat treatment, the strength increases to a limit, then begins to decrease. Aluminum-lithium alloys therefore come under the classification of precipitationstrengthening alloys. Since the size and distribution of the precipitates can be controlled by heat treating, these alloys are also referred to as heat-treatable. Since the addition of lithium to aluminum has also been shown to result in an alloy with unacceptable levels of ductility, a necessary engineering property for many aerospace applications, other elements such as copper, magnesium, and zirconium have been added to offset the loss in ductility. However, these alloy additions, particularly copper, increase the density. Consequently, schemes of alloy development have focused on balancing the various positive and negative attributes of the different elements, to arrive at a composition with suitable properties. Figure 1 is a schematic relating yield strength and
15 elongation, %
540
AI-Cu-Li-Mg-Zr 7075-T7651 modified 2020
AI-Li-Zr
10
AI-Mg-Li-Mn Al-Li-Mn 5
AI-Li-Cu-Zr AI-Li
100 (14.5)
200 (29)
300 (43.5)
400 (58)
500 (72.5)
yield strength, MPa (Ksi) Fig. 1. Strength–ductility relationships for various Al-Li-X alloys. The alloy 7075-T7651 is included as a reference point.
Aluminum alloys elongation for a variety of aluminum-lithium alloys. The alloy 7075–T7651, commonly used for such applications, is included as a reference point. In order to be competitive, the new alloys should have strength-ductility relations equal to or better than 7075. See LITHIUM. Powder metallurgy. As early as 1947, investigators showed interest in the development of highmodulus, high-temperature aluminum alloys. Elements such as iron, cobalt, and nickel have been added to aluminum in large quantities to produce alloys which have relatively stable structures at elevated temperatures and higher elastic moduli than conventional aluminum alloys. The titanium in the fan sections of gas turbines has to withstand operating temperatures up to 450◦F (235◦C). The substitution of a suitable aluminum alloy for titanium could result in a 31% reduction in weight for that part. In order to develop aluminum alloys with sufficient strength to withstand the operating temperatures of a gas turbine and remain thermally stable, researchers have considered alloy additions to aluminum which move slowly (have low diffusivities) in it, so that the characteristics of that structure do not change with time at the service temperature. Since the elements which have low diffusivities in aluminum also have very limited solubilities, processes other than conventional ingot casting must be utilized. An alternative for producing a candidate alloy is the powder metallurgy process, which is a rapid solidification process. Rapid solidification processing involves the transformation of finely divided liquid either in the form of droplets or thin sheets which become solid at high solidification rates (on the order of 104 K/s). Solidification occurs as heat is removed from the molten metal. The primary mechanisms of heat transfer can be divided into two broad categories—conductive and convective. Certain rapid-solidification-process forms such as ribbons and fibers are cooled primarily by conduction. The heat is transferred away by a substrate—the surface on which the liquid is cast. The substrate is generally made out of copper and can be either a rotating drum or a rotating wheel. On the other hand, another rapid-solidification-process form, powder, is cooled primarily by convection. The heat is transferred away from the fine liquid droplets in a gas stream. Regardless of the mechanisms of heat transfer, the resulting structures are more homogeneous and are on a finer scale than those formed by conventional ingot-casting techniques. The rapidly solidified particulates are then consolidated by forging or extrusion into a final shape. Even after consolidation, the various rapid-solidification-process forms are more homogeneous and have finer microstructures than a corresponding ingot-cast and worked product. Research centered on utilizing the natural benefit of rapid-solidification processing has produced metastable microstructures which decompose and precipitate during consolidation and working. Thus elements which have limited solubility and low diffusivity in aluminum (such as iron, cobalt, and
nickel—the transition metals), in conjunction with elements such as molybdenum or cerium, can be used to maximum advantage. However, the use of these alloying elements necessitates control of thermal history of the particulate both in its formation and upon subsequent processing. For example, in aluminum–transition metal alloys, the precipitates can form directly from the melt rather than from the solid. Consequently, the precipitates in these systems tend to be coarser than those in aluminum-lithium systems. As in the case of aluminum-lithium alloys, the yield strength is affected primarily by the interparticle spacing and the fineness of the precipitates which form. The ductility and fracture toughness are affected by macroparticles that concentrate stress, thereby reducing the toughness. Since ductility and fracture toughness are important properties, rapidsolidification-process alloys will be limited until a control of the volume fraction of coarse particles can be achieved. Rapid-solidification-process aluminum alloys containing iron and cerium appear to result in alloys which have refined particle-size distribution and improved high-temperature stability. These powder metallurgy alloys appear competitive with titanium up to 375◦F (191◦C). See METALLIC GLASSES; POWDER METALLURGY. Mechanical alloying. This process circumvents the limitations of conventional ingot casting. Blends of powders are mixed in a ball mill. A drum is mounted horizontally and half-filled with steel balls and blends of elemental metal powders. As the drum rotates, the balls drop on the metal powder. The degree of homogeneity of the powder mixture is determined by the size of the particles in the mixture. If the powder’s particles are too coarse, the different constituent elements in the blend will not interdiffuse during consolidation. Consequently, high rates of grinding must be used. In a high-energy mill the particles of metal are flattened, fractured, and rewelded. Every time two balls collide, they trap powder particles between them. The impact creates atomically clean fresh surfaces. When these clean surfaces meet, they reweld. Since such surfaces readily oxidize, the milling operation is generally conducted in an inert atmosphere. During the initial stages of the mechanical alloying process, the particles are layered composites of the constituent powders and the starting constituents are easily identifiable within the composite particle. The composition and shape of the individual particles are different. As alloying progresses, the particles are further fractured and rewelded. The intimate mixture of the starting constituents decreases the diffusion distance, and the particle tends to become more homogeneous. Metastable and stable phases are beginning to form in the individual particles. Diffusion of the elements is accelerated by the strain introduced during the milling. During the final stage of processing, the individual particles became similar in composition and similar to one another. Completion of the process occurs when the tendency to fracture is balanced by the tendency to
541
Alunite TABLE 3. Data for several loading fractions of silicon carbide (SiC) Tensile properties of SiC/2024 Al-T4 composites∗ 0 15
Vol. % SiC: Ultimate tensile strength, MPa (Ksi)† Elongation, GPa (Msi)§ Elongation to failure, %
Ultimate tensile strength, MPa (Ksi)† Elongation, GPa (Msi)§ Elongation to failure, % ∗
18 specimens.
NA‡
400–434 (58–64) 73 (10.6) 19–20
89–97 (13–14) 1–2
20
25
455–524 (66–76) 97–117 (14–17) 1–2
552–641 (80–93) 117–151 (17–22) 1–2
Tensile properties of SiC/2024 Al composites, as extruded ¶ 0 20
Vol. % SiC:
† Ksi
186 (27) 73 (10.6) 20–22
= thousand pounds per square inch.
331–372 (48–54) 97–117 (14–17) 1–2 ‡ Not
weld and the particle size distribution is constant. The relationship between the thickness and layers within the composite is schematically illustrated in Fig. 2. Aluminum–silicon carbide composites. A composite is a material in which two or more constituents are combined to result in a material which has properties different from those of either constituent. Typical composites are from materials in which one of the components has very high strength and modulus and the other has high ductility. Their properties generally follow a rule of mixtures. For example, if elastic modulus is the property of interest, the elastic modulus of the composite is approximately the weighted sum of the elastic moduli of the constituents. In one example, silicon carbide whiskers are mixed with aluminum alloy powder, compacted under pressure at elevated temperatures, and extruded or forged into a final product. The very highmodulus silicon carbide is incorporated in the duc-
available.
§ Msi = million pounds per square inch.
25 400–448 (58–65) 117–151 (17–22) 1–2 ¶ 20 specimens.
tile aluminum alloy matrix. The resulting properties depend on the volume fraction of silicon carbide, and to a large degree on the fabricating methods. The primary advantage of an aluminum–silicon carbide composite is the high elastic modulus and strength (Table 3). See ALUMINUM; COMPOSITE MATERIAL. Tom H. Sanders, Jr. Bibliography. American Society for Metals, Source Book on Selection and Fabrication of Aluminum Alloys, 1978; A. P Divecha, S. G. Fishman, and S. D. Karmarkar, Silicon carbide reinforced aluminum: A formable composite, J. Metals, 33:12–17, September 1981; P. S. Gilman and J. S. Benjamin, Mechanical alloying, Annu. Rev. Mater. Sci., 13: 279–300, 1983; F. King, Aluminum and Its Alloys, 1987; S. L. Langenbeck, Elevated temperature aluminum alloy development, Interim Technical Report for Period April–September 1981, U.S. Airforce Rep. LR 29977, October 1981; E. A. Starke, Jr., T. H. Sanders, Jr., and I. G. Palmer, New approaches to alloy development in the Al-Li system, J. Metals, 33:24–33, August 1981.
100
Alunite spacing, µm
542
10
1
0
50 processing time, min
100
Fig. 2. Effect of processing time on the spacing between the layers of mechanically alloyed powder. The smaller the spacing, the more homogeneous the powder particle.
A mineral of composition KAl3(SO4)2(OH)6. Alunite occurs in white to gray rhombohedral crystals or in fine-grained, compact masses. It has a hardness of 4 and a specific gravity of 2.6–2.8. Alunite is produced by sulfurous vapors on acid volcanic rocks and also by sulfated meteoric waters affecting aluminous rocks. The mineral is generally found associated with quartz and kaolinite. Alunite is used as a source of potash or for making alum. Alum has been manufactured from the wellknown alunite deposits at Tolfa, near Civita Vecchia, Italy, since the mid-fifteenth century. In the United States alunite is widespread in the West. Large deposits occur in Mineral and Hindsdale counties, Colorado, and in Piute County, Utah. See ALUM; FERTILIZER; POTASSIUM. Edward C. T. Chao
Alzheimer’s disease
Alzheimer’s disease A disease of the nervous system characterized by a progressive dementia that leads to profound impairment in cognition and behavior. Dementia occurs in a number of brain diseases where the impairment in cognitive abilities represents a decline from prior levels of function and interferes with the ability to perform routine daily activities (for example, balancing a checkbook or remembering appointments). Alzheimer’s disease is the most common form of dementia, affecting 5% of individuals over age 65. The onset of the dementia typically occurs in middle to late life, and the prevalence of the illness increases with advancing age to include 25–35% of individuals over age 85. See AGING. Clinical characteristics. Memory loss, including difficulty in remembering recent events and learning new information, is typically the earliest clinical feature of Alzheimer’s disease. As the illness progresses, memory of remote events and overlearned information (for example, date and place of birth) declines together with other cognitive abilities. Orientation to time (for example, date, day of week, or season) is often impaired, followed by a decline in orientation to personal information and immediate environment. Language deficits usually begin with difficulty in naming objects and in word-finding ability and gradually progress to a loss of all language skills. The ability to use visuospatial skills to draw simple figures and shapes is increasingly impaired. Declines in abstract reasoning, judgment, and problem solving interfere with the ability to perform mathematical calculations and to make appropriate decisions in social situations. In the later stages of Alzheimer’s disease, there is increasing loss of cognitive function to the point where the individual is bedridden and requires full-time assistance with basic living skills (for example, eating and bathing). See MEMORY. Although Alzheimer’s disease typically has a slowly progressing course, the rate of cognitive decline is variable among individuals. The average duration from the onset of symptoms to death is 8–10 years. In individuals reaching the end stage of dementia, the cause of death is often related to secondary medical conditions (for example, respiratory infections). Behavioral disturbances that can accompany Alzheimer’s disease include agitation, aggression, depressive mood, sleep disorder, and anxiety. Changes in personality, increased social withdrawal, and irritability are observed in some individuals. Delusions and hallucinations can occur. The onset of clinical symptoms can vary from middle to advanced age. Although a distinction between individuals with an early onset (65 years of age and younger) and with a late onset (older than 65) of dementia has been made, the same neuropathological disease process occurs in both groups. Neuropathological features. The major neuropathological features of Alzheimer’s disease include the presence of senile plaques, neurofibrillary tangles, and neuronal cell loss. Although the regional distri-
bution of brain pathology varies among individuals, the areas commonly affected include the association cortical and limbic regions. Senile plaques are extracellular and contain a core of an amyloid beta peptide that is derived from the cleavage of a much larger beta amyloid precursor protein. Although the amyloid precursor protein occurs normally in neural cell membranes, the reason why individuals with Alzheimer’s disease produce excessive amounts of the amyloid beta peptide is unknown. Neurofibrillary tangles represent intracellular pathology in which paired helical filaments accumulate within neurons. The basic component of neurofibrillary tangles is a hyperphosphorylated tau protein producing a highly insoluble cytoskeletal structure that is associated with cell death. Prominent loss of pyramidal neurons from cerebral cortical layers III and V has been observed in Alzheimer’s disease. Other microscopic features include neuropil threads, Hirano bodies, granulovascular degeneration, and inflammation around plaque formations. The brain regions of the medial temporal lobes are commonly affected early in the course of Alzheimer’s disease. In vivo neuroimaging studies that measure functional brain activity by assessing cerebral blood flow and metabolism have shown that the temporal, parietal, and frontal association cortices appear most consistently and severely affected in the early to middle stages of the illness. Alzheimer’s disease has been associated with dysfunction of several neurotransmitter systems. These neurochemical systems are responsible for the transfer of information between neurons via the synapse and are important for cognitive and behavioral functions. In Alzheimer’s disease, disruption of information transfer between neurons can occur with loss of presynaptic elements and reduction of the recognition sites at which the neurotransmitter acts. It is also thought that loss of or abnormalities in the receptor-operated ion channels and second messenger systems within cells may be related to neuronal dysfunction in Alzheimer’s disease. Deficits in cholinergic, serotonergic, noradrenergic, and peptidergic (for example, somatostatin) neurotransmitters have been demonstrated. Dysfunction of the cholinergic neurotransmitter system has been specifically implicated in the early occurrence of memory impairment in Alzheimer’s disease, and it has been a target in the development of potential therapeutic agents. The deficiency in choline acetyltransferase, an enzyme important for the synthesis of acetylcholine, has been associated with dementia severity. The cortical projection of the cholinergic system is especially concentrated in key structures such as the hippocampus in the medial temporal lobes, brain structures that are thought to be important in aspects of learning and memory. See ACETYLCHOLINE; NEUROBIOLOGY; NORADRENERGIC SYSTEM. Diagnosis. A definite diagnosis of Alzheimer’s disease is made only by direct examination of brain tissue obtained at autopsy or by biopsy to determine the presence of senile plaques and neurofibrillary tangles. A clinical evaluation, however, can provide
543
544
Amalgam a correct diagnosis in more than 80% of cases. The clinical diagnosis of Alzheimer’s disease requires a thorough evaluation to exclude all other medical, neurological, and psychiatric causes of the observed decline in memory and other cognitive abilities. This evaluation begins with a clinical history obtained from the patient and caregivers concerning the onset and course of symptoms and whether there are other medical conditions, head injury, excessive use of alcohol, or use of medications that could cause or exacerbate the cognitive decline. Neurological and psychiatric examinations are performed to characterize the cognitive deficits and to evaluate the presence of other disorders, such as Parkinson’s disease, Huntington’s disease, and cerebrovascular disease, that can cause dementia. The presence of behavioral disturbances such as depression and psychosis that can compromise cognitive function must also be evaluated. Laboratory tests are used to screen for medical conditions, including diabetes, renal dysfunction, hypothyroidism, infection, and vitamin deficiencies. A detailed neuropsychological evaluation is made to quantify and describe the cognitive deficits and to assist in staging the severity of dementia. Neuroimaging studies of the brain are performed to help rule out the presence of stroke, brain tumor, or other central nervous system abnormalities. Risk factors. Although the cause of Alzheimer’s disease is unknown, a number of factors that increase the risk of developing this form of dementia have been identified. Age is the most prominent risk factor, with the prevalence of the illness increasing twofold for each decade of life after age 60. Research in molecular genetics has shown that Alzheimer’s disease is etiologically heterogeneous. Gene mutations on several different chromosomes are associated with familial inherited forms of Alzheimer’s disease. Mutations of the amyloid precursor protein gene located on chromosome 21 have been linked to early-onset Alzheimer’s disease in a small number of families. Additionally, virtually all individuals with Down syndrome (where there is an extra copy of chromosome 21) develop the brain pathology of Alzheimer’s disease by the fourth decade of life. Gene defects on chromosomes 1 and 14 have also been implicated in familial early-onset forms of the illness, and they are thought to account for the vast majority of familially inherited cases of Alzheimer’s disease. In the more common sporadic, late-onset form of Alzheimer’s, specific allele types of the apolipoprotein gene on chromosome 19 act as a susceptibility factor that influences the age of onset of the disease. Treatment. A major strategy for the treatment of Alzheimer’s disease has focused on the relation between memory impairment and dysfunction of the acetylcholine neurotransmitter system. Tetrahydroaminoacridine (Tacrine) and donepezil hydrochloride (Aricept) were approved by the U.S. Food and Drug Administration for the treatment of Alzheimer’s disease. These medications serve to increase the amount of acetylcholine available in the synapse by preventing the breakdown of the neurotransmitter that naturally occurs with presence of the enzyme
acetylcholinesterase. Other treatment strategies to delay or diminish the progression of Alzheimer’s disease are being explored, including efforts to modulate multiple neurotransmitter systems, to reduce the accumlation of amyloid in the brain, and to diminish the inflammatory response. Behavioral and pharmacological interventions are also available to treat the specific behavioral disturbances that can occur in Alzheimer’s disease, such as depression, agitation, and sleep disorder. Gene Alexander Bibliography. M. M. Esiri and J. H. Morris (eds.), The Neuropathology of Dementia, Cambridge University Press, 1997; B. A. Lawlor (ed.), Behavioral Complications in Alzheimer’s Disease, American Psychiatric Press, 1995; R. D. Terry, R. Katzman, and K. C. Bick (eds.), Alzheimer’s Disease, Raven Press, 1994; W. Wasco and R. E. Tanzi (eds.), Molecular Mechanisms of Dementia, Humana Press, 1997.
Amalgam An alloy of mercury. Practically all metals will form alloys or amalgams with mercury, with the notable exception of iron. Amalgams are used as dental materials, in the concentration of gold and silver from their ores, and as electrodes in various industrial and laboratory electrolytic processes. Amalgams used in dental work require the following composition: silver, 65% minimum; copper, 6% maximum; zinc, 2% maximum; and tin, 25% minimum. These amalgams are prepared by the dentist as needed, and harden within 3–5 min, but may be shaped by carving for 15 min or so. See ALLOY; GOLD; MERCURY (ELEMENT); SILVER. E. Eugene Weaver
Amaranth An annual plant (seldom perennial) of the genus Amaranthus (family Amaranthaceae), distributed worldwide in warm and humid regions. Amaranths are botanically distinguished by their small chaffy flowers, arranged in dense, green or red, monoecious or dioecious inflorescences (spikes; see illus.), with zero to five perianth segments and two or three styles and stigmata, and by their dry membranous, indehiscent, one-seeded fruit. See FLOWER; FRUIT. Physiological, genetic, and nutritional studies have revealed their potential economic value. Of particular interest are high rate of productivity as a rapidly growing summer crop, the large amounts of protein in both seed and leaf with high lysine, the overall high nutritional value, and the water use efficiency for the C4 photosynthetic pathway. Amaranths are important in the culture, diet, and agricultural economy of the people of Mexico, Central and South America, Africa, and northern India. Genetic, ethnobotanical, and agronomic research has been undertaken to develop amaranths as an important food plant in modern agriculture. Classification and species relationships. Up to 200 or more species of Amaranthus have been
Amateur radio
3 cm
Flowers on branch of Amaranthus tricolor.
described, of which nearly 50 are distributed in the New World. Some of the more widespread weedy species are red amaranth or pigweed (A. hybridus), thorny amaranth (A. spinosus), tumbleweed (A. albus), and wild beet (A. hybridus ssp. hypochondriacus); Joseph’s coat (A. tricolor) and love-lies-bleeding (A. caudatus) are recognized for their ornamental spikes and foliage colors. Two sections of the genus are Amaranthotypus, with large terminal inflorescence, five tepals and stamens, and fruit opening circularly; and Blitopsis, with axillary cymes, often two to four tepals and stamens, and fruits opening irregularly. The cultivated grain species (Amaranthus cruentus, A. caudatus, and A. hypochondriacus = A. leucocarpus) with ornamental and crop forms belong to the section Amaranthotypus, whereas two important cultivated vegetable species (A. tricolor and A. lividus) belong to the section Blitopsis. Several Asian natives, such as A. tricolor, A. gangeticus, and A. melancholicus, have been grown as pot herbs or ornamentals but never as a grain crop. The cultivated grain species are largely paleseeded, whereas the ornamentals and weedy forms are black-seeded, presumably an important trait used by early peoples in distinguishing them. Interspecific hybridization, reported in various regions, is considered to be a primary source of evolutionary diversity. The grain crop forms and their weedy relatives include both 32- and 34-chromosome species; much of the preliminary biosystematic and cytogenetic research has not yet provided conclusive species relationships. However, the weedy and cultivated species of Amaranthotypus are separable into two distinct groups on the basis of taxonomic and protein-variation studies.
Domestication and cultivation. Amaranths are of ancient New World origin and were important as grains in the Aztec, Incan, and Mayan civilizations where wheat and rice were not grown. Archeological and ethnobotanical evidence suggests that the earliest domestication took place in South America at least 4000 years ago. About A.D. 500, grain amaranths were a widely cultivated food crop in Mexico. Some of the ornamental forms moved from the New World to Africa and Asia in the early 1800s, and may have spread from India to the East Indies, China, Manchuria, and Polynesia as a leaf vegetable. Ceremonial sacrifice and use in rituals drew the attention of the Spaniards, who suppressed the Aztec culture and amaranth cultivation. Amaranths are a widely grown vegetable in tropical regions, including southern India, West and Equatorial Africa, and New Guinea. Amaranths have been found to be the most popular vegetable in southern Benin, in West Africa. Varieties such as Chinese spinach, Ceylon spinach, and African spinach suggest use as a salad ingredient; however, somewhat high oxalic acid content make the leaves unpleasant for this purpose. Food value. Leaf amaranths are rich in provitamin A, vitamin C, iron, calcium, and protein, with lysine constituting as much as 5.9% of protein (equal to soymeal, and more than some of the best maize strains); glutamic acid constitutes as much as 10.8% of protein in A. edulis. Amaranth grains, with starch of very high quality and quantity, high protein (up to 17%), and high digestibility, rate better in nutritive value than all other cereals. They have been used in Mexico in modern times to make alegrias (cakes of popped seed), flour (of popped seed) used in the preparation of pinole, and paste used in the preparation of chuale; the black-seeded and variegated forms have been used to make amaranth milk (atole). In India, amaranth seed is used in alegrias and in confectioneries (flour) and as a cereal (popped seed). See CARYOPHYLLALES. Subodh K. Jain Bibliography. R. Bressani, Amaranth Newsletters, Arch. Latino Americanos de Nutricion, Guatemala, 1984–1997; G. J. H. Grubben, The Cultivation of Amaranth as a Tropical Leaf Vegetable, Roy. Trop. Inst. Commun. 67, 1976; S. K. Jain, P. A. Kulakow, and I. Peters, Genetics and breeding of grain amaranth: Some research issues and findings, Proceedings of the 3d Amaranth Conference, 1984, pp. 174–191, 1986; National Research Council, Amaranth: Modern Prospects for an Ancient Crop, 1984; R. M. Saunders and R. Becker, Amaranthus: A potential food and feed resource, Adv. Cereal Sci. Technol., 6:358– 396, 1984.
Amateur radio Two-way radio communications by individuals as a leisure-time activity. Amateur, or “ham,” radio is defined by international treaty as a “service of selftraining, intercommunications, and technical investigation carried on by amateurs; that is, by duly
545
546
Amateur radio authorized persons interested in radio technique solely with a personal aim and without pecuniary interest.” Activities. The government allows amateur operators many privileges because the hobby is partially based on service to the general public, and hams can be relied on to assist during emergencies. Groups of amateur operators meet annually to practice handling emergency communications in the field and to compete against other groups nationwide in performing certain emergency-related tasks. Amateur operators may set up warning and relief networks during the hurricane and tornado seasons, and handle communication when phone lines or cellular links are damaged by disasters. In addition to public service activities, amateurs enjoy many recreational activities, including DXing (where the objective is to contact amateurs in as many foreign countries as possible), contesting (where the amateurs compete for the maximum number of contacts in a given time span), and foxhunting (where the objective is to use radio skills to locate a hidden transmitter). Since high-frequency (HF) signals (below 30 MHz) are reflected from the Kennelly-Heaviside ionosphere layers, amateurs are commonly able to carry out international communication. Very high frequencies (VHF; 30–300 MHz) and ultrahigh frequencies (UHF; 300–3000 MHz) have been subject to exploration by amateurs at the leading edge of communications technology. Amateurs bounce signals off the Moon or ionized meteor trails, and communicate through amateur operator-built earth satellites called OSCAR (orbiting satellites carrying amateur radio). Equipment. For reliable, versatile, long-distance communications, some operators use large, directional beam antennas, but many have much simpler wire antennas strung between trees. Many operators have inexpensive equipment consisting of a wire antenna and just one transceiver (a radio that both receives and transmits) that functions on five or more frequency bands. See ANTENNA (ELECTROMAGNETISM). Unlike most radio services, amateurs may design and build their equipment as long as it meets FCC requirements for spectral purity and bandwidth. They may assemble it from kits or may purchase it ready to go. As a result of the proliferation of commercial equipment, operators generally build only accessory items or experiment with modifications to equipment they have purchased. A few operators design equipment for frequency bands where commercial choices are limited, or experiment with new devices originally designed for the commercial frequencies. Privileges. In the United States, the Federal Communications Commission (FCC) issues five classes of license, progressively allowing more frequencies and greater privileges. The first license level is the Novice Class, followed, in order of increasing privileges, by the Technician Class, General Class, Advanced Class, and Amateur Extra Class. In 1991 the FCC eliminated the Morse code requirement from the Technician Class license, which
shortly thereafter became the license of choice for most entry-level amateurs. Technician Class licensees have operating privileges on all amateur radio frequencies above 30 MHz. Novice Class licensees may use voice, Morse code, and digital communications on a limited number of frequencies that allow worldwide HF voice communications and short-distance VHF voice communications on frequencies where readily available handheld radios allow mobility and portability. Higher-class licensees have more frequency privileges and are permitted 1500 W peak envelope power output on most bands. Higher-class licensees may use any of the following modes: code (CW); amplitude-modulated (AM) and singlesideband (SSB) suppressed-carrier voice; frequency modulation (FM); radioteletype (RTTY), using either the Baudot, ASCII, or AMTOR (amateur teletype over radio) digital codes; facsimile; television; or pulse modulation. See AMPLITUDE-MODULATION RADIO; FREQUENCY-MODULATION RADIO; PULSE MODULATION; SINGLE SIDEBAND; TELETYPEWRITER; TELEVISION. Amateur operators with General Class and higher licenses can communicate on portions of or all of the following frequency bands, depending on their specific license level: 1.8–2.0 MHz 3.5–4.0 7.0–7.3 10.10–10.15 14.00–14.35 18.068–18.168 21.000–21.450 24.890–24.990 28.0–29.7 50.0–54.0 144.0–148.0 222–225 420–450 902–928
1240–1300 2300–2310 2390–2450 3300–3500 5650–5925 10,000–10,500 24,000–25,250 48,000–50,000 71,000–76,000 165,000–170,000 240,000–250,000 Others above 300,000 MHz
See RADIO SPECTRUM ALLOCATIONS. Requirements. Questions found on each license class examination involve basic radio theory, rules that the FCC requires amateur radio operators to abide by, and a knowledge of international Morse code (except the Technician Class license). In each license class exam the questions get progressively more difficult to cover topics that befit the accompanying privileges of that license. For example, because the General Class license allows RTTY, there are questions covering RTTY on that exam but not on the Novice Class exam. The Novice Class license requires the ability to send and receive Morse code at the rate of five words per minute (WPM); the General Class license requires 13 WPM; and the Amateur Extra, 20 WPM. Technical developments. The open policies toward amateur radio by the government have made pioneering possible in several areas. In 1923, for instance, amateurs led the way in the use of short
Amber waves after an amateur in Connecticut talked with a ham in France on the wavelength of 110 m (360 ft). The accepted theory had been that only wavelengths above 200 m (660 ft) were really useful for reliable communications. By the early 1930s communications between points on opposite sides of the globe had become commonplace on wavelengths of 80–20 m (264– 66 ft; frequencies of 3.5–14.4 MHz). Many hams then turned their attention to higher frequencies. Again amateurs shattered previous conceptions, especially those concerning so-called line-of-sight limitations on communications at 56 MHz and above. These pioneers discovered means of radio propagation which reached far beyond the usual horizon: reflections from the aurora borealis, from meteor trails, from sporadic E layers (patches of ionized particles about 70 mi or 110 km above the Earth), and from the Moon; bending of radio waves through layers of stable air; and a phenomenon called transequatorial scatter, by which stations on one side of the Equator may communicate with stations on the other side over distances of more than 1000 mi (1600 km) at times when such communication would otherwise be considered impossible. This mode and several of the others were subjects of an International Geophysical Year study undertaken by amateurs through the American Radio Relay League. See RADIO-WAVE PROPAGATION; TROPOSPHERIC SCATTER. Other technical accomplishments and experimentation by amateurs include computer communications multiplexed on common radio channels, or packet radio; computer station control and interfacing as a result of the popularity and availability of personal computers; and spread spectrum techniques, enabling hundreds of simultaneous conversations to be carried on one frequency band without any mutual interference. Amateurs have installed relay repeater stations on mountaintops in many areas and use them to extend the range of communication on the VHF and UHF bands from tens to hundreds of miles. Under the direction of the Radio Amateur Satellite Corporation (AMSAT), amateurs have built and launched as secondary payloads on scheduled space shots a series of highly advanced satellites, some into low earth orbit and others into high elliptical orbits, containing refirable hydrazine engines, master on-board computers, and high-power VHFUHF transponders. These satellites act as relay stations, again increasing the range of possible communications on the frequencies of the transponders. See SPREAD SPECTRUM COMMUNICATION. Associations. The developments in the hobby over the years required a clearinghouse, an information exchange, which is embodied in the American Radio Relay League (ARRL), an association of radio amateurs in North America founded in 1914. It publishes a monthly magazine, a handbook, and other publications and operator aids covering various aspects of the hobby. ARRL also serves to represent amateurs in government regulatory matters; presents new technical developments and sponsors operating contests and other activities; and organizes networks to han-
dle messages and coordinates emergency communications training. The League is also Secretariat for the International Amateur Radio Union, composed of over 125 similar national radio societies around the world. See RADIO. Steve Mansfield Bibliography. American Radio Relay League: Advanced Class License Manual, Extra Class License Manual, The FCC Rule Book, Radio Amateur’s Handbook, General Class License Manual,Now You’re Talking, Radio Amateur’s Handbook, all annually.
Amber Most commonly, a generic name for all fossil resins, although it has been restricted by some to refer only to succinite, the mineralogical species of fossil resin making up most of the Baltic Coast deposits. Resins generally are complex mixtures of mono-, sesqui-, di-, and triterpenoids; however, some resins contain aromatic phenols or are even predominantly composed of these compounds. Resins are synthesized in appreciable quantity by about 10% of present-day plant families. Among the plants, primarily trees, that produce copious amounts of resin that may fossilize to become amber, about two-thirds are tropical or subtropical. Members of the families Pinaceae, Araucariaceae, and Taxodiaceae are the most prominent copious producers among the conifers, and Leguminosae, Dipterocarpaceae, Burseraceae, and Hamamelidaceae among the angiosperms. See RESIN. Fossilization. Under natural forest conditions resins harden with varying degrees of rapidity, and can become fossilized if they have the requisite chemical composition for polymerization, are sufficiently stable to withstand both oxidative and microbial degradation, and appropriate depositional conditions are available. The fossilization process is still not completely understood but appears primarily to involve polymerization, with some modification of the original constituents, especially due to the temperature during the maturation of the resin in the sediments. Occurrence. Although ambers occur throughout the world in deposits from Carboniferous to Pleistocene in age, they have been reported most commonly from Cretaceous and Tertiary strata and often are associated with coal or lignites. Amber may contain beautifully preserved insects, spiders, flowers, leaves, and even small animals (Fig. 1). The most extensively studied deposits are those from the Baltic Coast, Alaska, Canada, North American Atlantic coast, southeast Asia, Dominican Republic, and Mexico (Fig. 2). Uses. Since the earliest stages of human social development, ambers have had esthetic appeal, have been used to ward off evil powers, and have been thought to cure certain illnesses. The attribution of these special powers to amber may result partially from the negative electrical properties exhibited when ambers are rubbed. Thales recognized
547
548
Amber as during Roman and Greek periods. When used for jewelry, it usually is transparent yellow, reddishbrown, or “amber” color. Translucent or semitranslucent amber is used for pipe stems, decorating small boxes, and other ornamental purposes. Specific gravity varies from 1.05 to 1.10, and hardness from 1 to 3 on Mohs scale. Large amounts of amber occur in some coal deposits, which are being modified for some industrial uses, such as jet fuel. Furthermore, fossil resins, probably from members of the Dipterocarpaceae, have made marked contributions to petroleum in numerous southeast Asian sites. Chemistry and origin. At one time, chemical studies of amber were mineralogically oriented because the purpose was to describe and classify amber as a semiprecious gem. However, phytochemical studies comparing fossil and present-day resins, employing such techniques as infrared spectroscopy, x-ray diffraction, nuclear magnetic resonance, and pyrolysis gas chromatography–mass spectrometry are providing information regarding the botanical origins of ambers. For example, amber from Chiapas, Mexico, and from the Dominican Republic has been shown to be derived from the leguminous genus Hymenaea; that from Sumatra from Shorea (Dipterocarpaceae); that from Ecuador from members of the Burseraceae; that from Australia and New Zealand
Fig. 1. Pseudosphegina carpenteri (Diptera), syrphid fly in Baltic amber, Oligocene. (Courtesy of F. M. Carpenter)
these properties and, in fact, the term electricity evolved from electron, the Greek name for amber. See ELECTRIC CHARGE; ELECTRICITY. Carvings and beads of Baltic amber found in caves and tombs indicate its significance in human activities since the Stone Age. Amber provided a distinctive and imperishable barter item for trade routes that crossed Europe from the Baltic to the Adriatic and Black seas to the Bronze and Iron ages as well 120°
150°
180°
150°
120°
90°
4
60°
30°
0°
30°
60°
5
C
4
60°
60° U
T
U
30° 5
C
T
5 C 1
U T 4
C 4 5 2 C 4
9
8
U
1 5
3 4
3 3
30°
T7
7
U 0°
1 7 5 57 5 C2 6T C 7 4C C 7
T 7
T
0°
T
8
7 9
C
30°
9 8
8
Mercator projection 0 0
150°
180°
1 Carboniferous 2 Triassic 3 Early Cretaceous 4 Late Cretaceous 5 Eocene 6 Oligocene
2000 mi at Equator 2000 km
150°
120°
30°
Key:
60°
120°
9
90°
60°
30°
0°
7 Miocene 8 Pliocene 9 Pleistocene T Undetermined Tertiary U Unknown C Undetermined Cretaceous
30°
60°
60°
Fig. 2. Important amber deposits of the world. (After J. H. Langenheim, Amber: A botanical inquiry, Science, 163:1157–1169, 1969)
Amblygonite from Agathis (Araucariaceae); and that from Alaska probably from genera, such as Metasequoia, in the Taxodiaceae. The origin of the amber from the extensive Baltic deposits remains a mystery as the chemistry of the resin is more closely related to Agathis than to Pinus (Pinaceae), which early was considered to be the source from parts of the plant found in the amber. Then Agathis was considered because it has many chemical similarities to Baltic amber; however, it lacks succinic acid, a characteristic component, as well as plant remains of it in the amber. Recently amber from a Canadian Arctic site that contains succinic acid has been found within cone scales of Pseudostrobus (Pinaceae). As with Agathis, no parts of Pseudostrobus have been found in the amber. Although some chemical aspects differ from Baltic amber, the presence of succinic acid in the Pseudostrobus amber may provide a clue for future searches regarding the enigmatic source of Baltic amber. The predominantly tropical or subtropical occurrence of amber-producing plants through geologic time has led to evolutionary studies of the natural purpose of resins and their possible defensive role for trees against injury and disease inflicted by the high diversity of insects and fungi in tropical environments. See ARCHEOLOGICAL CHEMISTRY. Jean H. Langenheim Bibliography. D. A. Grimaldi, Amber: Window to the Past, Harry Abrams, New York, 1996; J. H. Langenheim, Biology of amber-producing trees: Focus on case studies of Hymenaea and Agathis, in K. B. Anderson and J. C. Crelling (eds.), Amber, Retinite and Fossil Resins, ACS Ser., no. 617, 1995; G. O. Poinar, Jr., and R. Poinar, The Amber Forest, Princeton University Press, 1999; P. C. Rice, Amber: The Golden Gem of the Ages, Kosciusko Foundation, New York, 1987.
Ambergris A fatty substance formed in the intestinal tract of the sperm whale (Physeter catodon). There is a question as to whether the origin of ambergris is normal or pathological, but it does serve as protection from the horny indigestible portions of the squid and cuttlefish that constitute much of the whale’s diet. Ambergris contains acids, alkaloids, and a fatty substance, its main constituent, called ambrein. Although fresh ambergris is soft and black and has an offensive odor, it hardens into pleasantly fragrant gray or yellow masses when exposed to the air, sun, and sea. Being lighter than water, it is found in lumps, weighing from 1/2 oz (15 g) to over 100 lb (45 kg) floating on tropical seas or cast up on the shores. It is also gathered directly from the abdomens of dead or captured whales when they are slaughtered. Collecting grounds for ambergris are principally on the shores of China, Japan, Africa, the Americas, tropical islands, and the Bahamas. Known and highly prized since the earliest times, ambergris was used chiefly for medicinal purposes.
In modern times this rare commodity is valued in the manufacture of perfumes. The ambergris is ground and used in the form of a tincture, dissolved in a dilute solution of alcohol, which when added to perfume acts as a fixative, increasing the duration of the fragrance while adding its own sweet, earthy scent. Sybil P. Parker
Amblygonite A lithium aluminum phosphate mineral of basic formula LiAl(PO4)(F). The structure of amblygonite consists of phosphate (PO4) groups of tetrahedra and AlO6 groups of octahedra. Each PO4 tetrahedron is connected to an AlO6 octahedron. Cornersharing octahedra form zig-zag chains along the b axis. Lithium (Li) is in fivefold coordination, and lies between the PO4 tetrahedra and nearest AlO6 octahedra. The dominant substitution in this mineral structure is hydroxyl (OH) for fluorine (F). This substitution gives rise to the amblygonite-montebrasite [LiAlPO4(OH)] solid solution series. When OH is greater than F, the mineral is known as montebrasite. Appreciable amounts of sodium substitute for lithium in the five-coordinated polyhedra. The sodium-rich varieties such as natromontebrasite and hedronite are rare. Amblygonite crystallizes in the triclinic system. Its color is commonly white or gray with tints of blue, green, and yellow. Amblygonite is transparent to translucent and has a vitreous to pearly luster. It often exhibits polysynthetic twinning on both a hand specimen and a microscopic scale. Two cleavage directions are prominent. Cleavage fragments of amblygonite-montebrasite may be confused with potassium feldspar, but they are distinguished by their higher density (specific gravity = 3.0–3.1) and cleavage angles that are not at 90◦. Members of the amblygonite-montebrasite series typically occur as coarse, cleavable nodules in the inner zones and quartz cores of zoned lithium-rich, granitic pegmatites. Quartz, spodumene, petalite, and feldspar are the most common minerals associated with amblygonite. Minerals such as lepidolite, beryl, cassiterite, tourmaline, and niobium-tantalum oxides are commonly associated with amblygonitebearing mineral assemblages in smaller abundances. These large masses of amblygonite-montebrasite are magmatic in origin and generally contain 4–7 wt % fluorine. Clear, equant crystals are rare and occur in vugs in pegmatites. Crystals associated with vugs and altered (secondary) magmatic amblygonite contain 0.3–4 wt % fluorine. A single mass of amblygonite with dimensions of 7.62 m × 2.44 m × 1.83 m (25 ft × 8 ft × 6 ft) has been documented in the Hugo granitic pegmatite in the Black Hills of South Dakota. Larger masses have been reported. The best-known occurrences of amblygonite are in Montebras, France; the Black Hills of South Dakota; the White Picacho District in Arizona; pegmatite districts in Maine; the Tanco pegmatite in Manitoba, Canada; and Portland, Connecticut. While
549
550
Amblypygi amblygonite has been mined as an ore of lithium, it is not a major ore. See PHOSPHATE MINERALS. Charles K. Shearer Bibliography. C. Klein and C. S. Hurlbut, Jr., Manual of Mineralogy, 21st ed., 1999; P. B. Moore, Pegmatite Minerals of P (V) and B (III): Granitic Pegmatities in Science and Industry, 1982.
Amblypygi The tailless whip scorpions or whip spiders, an order of the class Arachnida. There are about 80 species in the tropics and subtropics. They are flattened, red to brown, and range from 5 to 45 mm (0.2 to 1.8 in.) in body length. All are nocturnal, hiding under bark or stones, among leaves, or in caves during the day. The cephalothorax (prosoma) and segmented abdomen (opisthosoma) are narrowly attached. There is no tail. The cephalothoracic shield (carapace) has three simple eyes on each side and a pair of eyes anterior. The chelicerae are like the jaws of a spider. The raptorial pedipalps are strong and have long, hard spines used to surround and crush insect prey. The long, whiplike first legs are sense organs used to locate prey. There are booklungs on the underside of the abdomen. There are no venom or repellent glands. During courtship the male uses his long first legs to stroke the female; he may tremble while stroking. After some time he deposits a spermatophore (a package of sperm) on the ground. He pulls or guides the female over the spermatophore by using his chelicerae or first legs. The female picks up the spermatophore with her gonopore. Up to 60 eggs are laid in a flexible sac secreted by the female reproductive organs; the sac remains attached to the underside of the female’s abdomen. On hatching, the young climb up on their mother’s abdomen and are carried until their next molt. See ARACHNIDA. H. W. Levi
Ameba Any protozoon moving by means of protoplasmic flow. In their entirety, the ameboid protozoa include naked amebas, those enclosed within a shell or test, as well as more highly developed representatives such as the heliozoians, radiolarians, and foraminiferans. Ameboid movement is accomplished by pseudopods—cellular extensions which channel the flow of protoplasm. Pseudopods take varied forms and help distinguish among the different groups. A lobe-shaped extension or lobopod is perhaps the simplest type of pseudopod. The shapelessness and plasticity of these locomotory organelles impart an asymmetric, continually changing aspect to the organism. Other, more developed, representatives have pseudopodial extensions containing fibrous supporting elements (axopods) or forming an
extensive network of anastomosing channels (reticulopods). Though involved in locomotion, these organelles are also functional in phagocytosis—the trapping and ingesting of food organisms (usually bacteria, algae, or other protozoa) or detritus. See FORAMINIFERIDA; HELIOZOIA; PHAGOCYTOSIS; RADIOLARIA. Amebas are found in a variety of habitats, including fresh-water and marine environments, soil, and as symbionts and parasites in body cavities and tissues of vertebrates and invertebrates. Because of their manner of locomotion, amebas typically occur on surfaces, such as the bottom of a pond, on submerged vegetation, or floating debris. In soil, they are a significant component of the microfauna, feeding extensively on bacteria and small fungi. Amebas in marine habitats may be found as planktonic forms adapted for floating at the surface (having oil droplets to increase bouyancy and projections to increase surface area), where they feed upon bacteria, algae, and other protozoa. Several species of amebas may be found in the human intestinal tract as harmless commensals (for example, Entamoeba coli) or as important parasites responsible for amebic dysentery (E. histolytica). Amebas range from small soil organisms, such as Acanthamoeba (20 micrometers), to the large freshwater forms Amoeba proteus (600 µm; see illus.) and Pelomyxa (1 mm, or more). Some types, such as Amoeba, are uninucleate; others are multinucleate. Reproduction is by mitosis with nuclear division preceding cytoplasmic division to produce two daughters. Multinucleate forms have more unusual patterns of division, since nuclear division is not immediately nor necessarily followed by cytoplasmic division. Transformation of the actively feeding ameba into a dormant cyst occurs in many species, particularly those found in soil or as symbionts. The resting stages allow survival over periods of desiccation, food scarcity, or transmission between hosts. See REPRODUCTION (ANIMAL). Much attention has focused on the mechanism of ameboid movement in Amoeba proteus, a favorite research species because of its large size and ease of maintenance in the laboratory. Under the light microscope, the organism is seen to contain a granular endoplasmic core in a sol (fluid) state, enclosed by a gel-like ectoplasm. Endoplasm flows forward as the organism progresses; as it reaches the front end of the advancing pseudopod, the endoplasm is deflected sideways and rearward, where it converts to ectoplasmic gel. This pattern has been termed fountain streaming. At the posterior end of the organism, the reverse occurs, with ectoplasm reverting to a forward-flowing endoplasm. The ameba resembles an ectoplasmic cylinder enclosing a fluid endoplasmic core. It is generally accepted that movement is brought about by the contraction of protoplasm at the anterior end of the pseudopod, pulling the ameba along the substrate. The contraction involves cytoplasmic microfilaments, adenosine triphosphate (ATP) as the energy source, and divalent
Americium providing support for this evolutionary pathway. Among certain specialized ameboid protozoa, production of flagellated gametes prior to sexual reproduction may occur. See AMOEBIDA; PROTOZOA. Frederick L. Schuster
Americium A chemical element, symbol Am, atomic number 95. The isotope 241Am is an alpha emitter with a halflife of 433 years. Other isotopes of americium range in mass from 232 to 247, but only the isotopes of mass 241 and 243 are important. The isotope 241Am is routinely separated from “old” plutonium and sold for a variety of industrial uses, such as 59-keV gamma sources and as a component in neutron sources. The longer-lived 243Am (half-life 7400 years) is a precursor in 244Cm production. 1 1 H 3 Li 11 Na 19 K 37 Rb 55 Cs 87 Fr
2 4 Be 12 Mg 20 Ca 38 Sr 56 Ba 88 Ra
3 21 Sc 39 Y 71 Lu 103 Lr
4 22 Ti 40 Zr 72 Hf 104 Rf
lanthanide series actinide series
Phase-contrast photomicrograph of Amoeba proteus, a large fresh-water ameba. The organism is seen moving by means of a single lobose pseudopod.
cations, for example, Ca2+ and Mg2+. See CELL MOTILITY. Cytoplasm of Amoeba contains, besides food vacuoles, a contractile vacuole functioning for osmoregulation—a microkidney, which expells water. Energy for this, as well as other functions, is derived from numerous mitochondria present in the cytoplasm, which may also contain crystals and endosymbiotic bacteria. See MITOCHONDRIA; OSMOREGULATORY MECHANISMS; VACUOLE. Ameboid protozoa are believed to have evolved from flagellated protozoan ancestors through loss of the flagellar apparatus. Some ameboid organisms, the ameboflagellates, can change from creeping amebas to swimming flagellates during their life cycle,
5 23 V 41 Nb 73 Ta 105 Db
6 24 Cr 42 Mo 74 W 106 Sg
7 25 Mn 43 Tc 75 Re 107 Bh
8 26 Fe 44 Ru 76 Os 108 Hs
9 27 Co 45 Rh 77 Ir 109 Mt
10 28 Ni 46 Pd 78 Pt 110 Ds
11 29 Cu 47 Ag 79 Au 111 Rg
12 30 Zn 48 Cd 80 Hg 112
13 5 B 13 Al 31 Ga 49 In 81 Tl 113
14 6 C 14 Si 32 Ge 50 Sn 82 Pb
15 16 7 8 N O 15 16 P S 33 34 As Se 51 52 Sb Te 83 84 Bi Po
57 58 59 60 61 62 63 64 65 La Ce Pr Nd Pm Sm Eu Gd Tb
66 67 Dy Ho
89 Ac
98 Cf
90 Th
91 Pa
92 93 94 95 96 97 U Np Pu Am Cm Bk
18 2 17 He 9 10 F Ne 17 18 Cl Ar 35 36 Br Kr 53 54 I Xe 85 86 At Rn
68 69 70 Er Tm Yb
99 100 101 102 Es Fm Md No
In its most prominent aqueous oxidation state, 3+, americium closely resembles the tripositive rare earths. The formal analogy to the rare earths is also marked in anhydrous compounds of both tripositive and tetrapositive americium. Americium is different in that it is possible to oxidize Am3+ to both the 5+ and 6+ states. Americium metal has a vapor pressure markedly higher than that of its neighboring elements and can be purified by distillation. The metal is nonmagnetic and superconducting at 0.79 K. Under high pressure the metal has been compressed to 80% of its roomtemperature volume and displays the α-uranium structure. Americium-241 is used as an alpha-particle source in ionization-type smoke detectors. See ACTINIDE ELEMENTS; ALPHA PARTICLES; BERKELIUM; CURIUM; FIRE DETECTOR; NUCLEAR REACTION; PERIODIC TABLE; TRANSURANIUM ELEMENTS. R. A. Penneman Bibliography. N. M. Edelstein, J. O. Navratil, and W. W. Schulz, Americium and Curium Chemistry and Technology, 1985; S. Hofmann, On Beyond Uranium: Journey to the End of the Periodic Table, 2002; J. J. Katz, G. T. Seaborg, and L. R. Morss (eds.), The Chemistry of the Actinide Elements, 2 vols., 2d ed., 1986; W. W. Schulz, The Chemistry of Americium, ERDA Crit. Rev. Ser. Rep. TID-26971, 1976.
551
552
Amethyst
Amethyst The transparent purple to violet variety of the mineral quartz. Although quartz is perhaps the commonest gem mineral known, amethyst is rare in the deep colors that characterize fine quality. Amethyst is usually colored unevenly and is often heated slightly in an effort to distribute the color more evenly. Heating at higher temperatures usually changes it to yellow or brown (rarely green), and further heating removes all color. The principal sources are Brazil, Arizona, Uruguay, and Russia. Amethyst is often cut in step or brilliant shapes, and drilled or carved for beads. Carvings are made both from transparent and nontransparent material. See GEM; QUARTZ. Richard T. Liddicoat, Jr.
Formation and properties. Commercial preparation of amides involves thermal dehydration of ammonium salts of carboxylic acids. Thus, slow pyrolysis of ammonium acetate, CH3COO−NH4+, forms water and acetamide. N,N-dimethylacetamide may be similarly prepared from dimethylammonium acetate: +
CH3 COO− NH2 (CH3 )2
Acid anhydrides react with ammonia or with primary or secondary amines to form amides. The preparation of acetanilide [reaction (1)] is an example. NH2 + (CH3CO)O
NHCOCH3 + CH3COOH
(1)
Amide A derivative of a carboxylic acid with general formula (1), where R is hydrogen or an alkyl or O NH2
C
R
The reaction of amines with acid chlorides in the presence of aqueous sodium hydroxide (Schotten– Baumann reaction) is often used [reaction (2)]. COCI + RNH2 + NAOH
(1) aryl radical. Amides are divided into subclasses, depending on the number of substituents on nitrogen. The simple, or primary, amides are considered to be derivatives formed by replacement of the carboxylic hydroxyl group by the amino group, NH2. They are named by dropping the “-ic acid” or “-oic acid” from the name of the parent carboxylic acid and replacing it with the suffix “amide,” as shown in examples (2–5). In the secondary and tertiary amides, one or
NH2
Formamide
CO(NH2)2 + CH3COOH CH3CONH2 + H2NCOOH
NH2
CH3C
Acetamide
(2)
(4)
CH3 CH2 CH2 CH2 CONH2
N(CH3)2
N,N -Dimethylpropanamide
(3)
CH3 CH2 CH2 CH2 C N + H2 O −→
O CH3CH2C
+ NH3
The partial hydrolysis of nitriles also affords a convenient synthesis of a variety of amides [reaction (4)]. The reaction is catalyzed by both acid and base.
NHCH2CH2
N -Ethylbenzamide
CO2
(3)
O C
(2)
When urea, the diamide of carbonic acid, is heated with a carboxylic acid, a new amide is formed by an exchange reaction. The other product of the reaction, carbamic acid, decomposes to permit the reaction to go to completion [reaction (3)].
O
O HC
CONHR + NaCI + H2O
(5)
both hydrogens are replaced by other groups. The presence of such groups is designated by the prefix capital N (for nitrogen), as shown in the examples. Except for formamide, all simple amides are relatively low-melting solids, stable, and weakly acidic. They are strongly associated through hydrogen bonding, and hence soluble in hydroxylic solvents, such as alcohol. Because of ease of formation and sharp melting points, amides are frequently used for the identification of organic acids and, conversely, for the identification of amines.
(4)
Uses. Amides are important chemical intermediates since they can be hydrolyzed to acids, dehydrated to nitriles, and degraded to amines containing one less carbon atom by the Hofmann reaction. In pharmacology, acetaminophen (6) is a popular analgesic. N,N-Dimethylformamide (DMF; 7) is a useful aprotic, highly polar solvent. O CH3CONH
OH
(6)
HCN(CH3)2
(7)
The amides of the straight-chain fatty acids having 12, 14, 16, or 18 carbon atoms are especially useful industrially because of the variety of properties
Amine that can be obtained by substitution on the nitrogen atoms. Such products find extensive use as waterproofing agents, lubricant additives, detergents, emulsifiers, and wetting agents. Amides are very prevalent in nature, since all peptides and proteins are polymers of the natural αamino acids, as represented by (8). R O NHCHC
n
(8) See PEPTIDE; PROTEIN. Nylon is the generic term for any synthetic, longchain polymer with an amide linkage as a recurring, integral part of the chain. Nylon-6 (9) is one of the most useful. [
NH(CH2)5CO
]n
(9) See ACID ANHYDRIDE; ACID HALIDE; AMINE; CARBOXYLIC ACID; NITRILE; POLYAMIDE RESINS. Paul E. Fanta Bibliography. R. J. Fessenden and J. S. Fessenden, Organic Chemistry, 6th ed., 1998.
Amiiformes An order of actinopterygian fishes in the subclass Neopterygii. Amiiformes plus several related fossil orders comprise the Halecomorphi, a welldeveloped group known from the middle Mesozoic. The Amiiformes comprise several families that persisted into the Cenozoic era; but only one, the Amiidae, survived to the present and it consists of a single species, Amia calva, the bowfin. The fossil amiids are from freshwater deposits, the oldest of which are of Jurassic age. The other halecomorphs were marine species. Amia calva (see illustration) is characterized by an abbreviated heterocercal tail; fusiform body; large mouth equipped with sharp teeth; an elongate dorsal fin, with each ray supported by a single pterygiophore; scales with a ganoin (enamel-like) surface but thin and overlapping; no spiracles; a vascular swim bladder that can act as a lung; and a large gular
plate. The maximum length of the bowfin is about 90 cm (35 in.). Its usual habitat is still waters with an abundance of rooted vegetation in lowlands of eastern North America from the St. Lawrence River to the Gulf Slope and from the Mississippi Basin to the Atlantic Slope, avoiding Appalachian streams. See ACTINOPTERYGII; OSTEICHTHYES; SWIM BLADDER; TELEOSTEI. Herbert Boschung Bibliography. L. Grande and W. E. Bemis, A comprehensive phylogenetic study of amiid fishes (Amiidae) based on comparative skeletal anatomy, J. Vert. Paleontol., Spec. Mem. 4 (suppl. to vol. 18), 1998; J. G. Maisey, Santana Fossils: An Illustrated Atlas, T. F. H. Publications, Neptune City, NJ, 1991; J. S. Nelson, Fishes of the World, 4th ed., Wiley, New York, 2006.
Amine A member of a group of organic compounds which can be considered as derived from ammonia by replacement of one or more hydrogens by organic radicals. Generally amines are bases of widely varying strengths, but a few which are acidic are known. Amines constitute one of the most important classes of organic compounds. The lone pair of electrons on the amine nitrogen enables amines to participate in a large variety of reactions as a base or a nucleophile. Amines play prominent roles in biochemical systems; they are widely distributed in nature in the form of amino acids, alkaloids, and vitamins. Many complex amines have pronounced physiological activity, for example, epinephrine (adrenaline), thiamin or vitamin B1, and Novocaine. The odor of decaying fish is due to simple amines produced by bacterial action. Amines are used to manufacture many medicinal chemicals, such as sulfa drugs and anesthetics. The important fiber nylon is an amine derivative. Amines are classified according to the number of hydrogens of ammonia which are replaced by radicals. Replacement of one hydrogen results in a primary amine (RNH2), replacement of two hydrogens results in a secondary amine (R2NH), and replacement of all three hydrogens results in a tertiary amine (R3N). The substituent groups (R) may be alkyl, aryl, or aralkyl. In another group of amines the nitrogen forms part of a ring (heterocyclic amines). Examples of such compounds are nicotine (1), which is obtained commercially from tobacco for use as an insecticide, and serotonin (2), which plays N
CH3 N
(1) CH2CH2NH2
HO
Amia calva, the bowfin. (Courtesy of Frank Teigler, Hippocampus Bildarchiv)
N H
(2)
553
554
Amine a key role as a chemical mediator in the central nervous system. Many aromatic and heterocyclic amines are known by trivial names, and derivatives are named as substitution products of the parent amine. Thus, C6H5NH2 is aniline and C6H5NHC2H5 is N-ethylaniline. For a discussion of definitive rules for naming amines see ORGANIC CHEMISTRY
According to the Br¨ onsted-Lowry theory of acids and bases, amines are basic because they accept protons from acids. In water the equilibrium shown in reaction (1) lies predominantly to the left. The ex+ − RNH2 + H2 O − − RNH3 + OH
[RNH3 ]+ [OH]− RNH2
RNH2 + CHCl3 + 3OH− −→
(1)
tent to which the amine is successful in picking up a proton from water is given by Eq. (2), where the Kb =
character. Primary amines give sulfonamides that are soluble in basic solutions; secondary amines give insoluble derivatives; and tertiary amines with no replaceable hydrogen do not react with the reagent. In general, the sulfonamides are solids and are useful for identification of the amines. Carbylamines (isocyanides) possess a very unpleasant, nauseating odor and are formed by the reaction of any primary amine with chloroform in basic solution [reaction (5)].
(2)
quantities in brackets signify concentrations of the species given. For short-chain aliphatic amines the basic ionization constant Kb lies near 10−4; for aromatic amines Kb < 10−9; for ammonia Kb = 1.8 × 10−5. Stable salts suitable for the identification of amines are in general formed only with strong acids, such as hydrochloric, sulfuric, oxalic, chloroplatinic, or picric. Reactions and identification. Several test-tube reactions for recognition and characterization of amines are known: the Schotten-Baumann reaction, the Hinsberg test, the carbylamine reaction, and the action of nitrous acid. The Schotten-Baumann reaction involves treatment of an amine with benzoyl chloride in basic solution. It serves to distinguish tertiary amines from primary and secondary amines by the formation of substituted benzamides from the primary and secondary amines [reactions (3) and (4)]. Tertiary RNH2 + C6 H5 COCl + OH− −→ C6 H5 CONHR + Cl− + H2 O (3) R2 NH + C6 H5 COCl + OH− −→ C6 H5 CONR2 + Cl− + H2 O (4)
amines do not react with benzoyl chloride. The substituted benzamides are generally insoluble in water, solid, and easily purified, and they have characteristic melting points. Thus they serve to identify the amines. The more reactive acylating agents, acetic anhydride and acetyl chloride, give substituted acetamides without added base. This reaction gives the same type of information as the Schotten-Baumann reaction. Closely related to the Schotten-Baumann reaction is the Hinsberg test. This has the added advantage of distinguishing between primary and secondary amines. It involves reaction of an amine with benzenesulfonyl chloride in alkaline solution. Both it and the Schotten-Baumann test are applicable to both aliphatic and aromatic amines with the exception of those amines which are substantially nonbasic in
R
N
C + 3Cl− + 3H2 O (5)
Reaction with nitrous acid serves as a further method for distinguishing between various classes of amines. Primary aliphatic amines evolve nitrogen, whereas primary aromatic amines give diazonium compounds, which may be recognized by dye formation on coupling with a suitable second component. Secondary amines of both series give nitrosamines, generally as yellow oils. Tertiary aliphatic amines do not react with nitrous acid, and mixed aliphatic aromatic tertiary amines undergo nuclear nitrosation. Nitrosamines have been identified as potent carcinogenic substances. See DIAZOTIZATION. In the infrared absorption spectrum, amines exhibit characteristic bands due to N-H stretching and bending, as well as C-N stretching vibrations. In the proton magnetic resonance spectrum, the appearance of the proton on nitrogen is complicated by the rate of exchange and the electrical quadrupole of the 14N nucleus. In alkaline solutions tertiary amines are oxidized by hydrogen peroxide to amine oxides which, although still basic, are not strong bases as are the quaternary ammonium types [reaction (6)]. R3 N + H2 O2 −→ R3 NO + H2 O
(6)
Aromatic amines undergo halogenation and sulfonation on the ring. However, because of their susceptibility to oxidation, nitration cannot be accomplished without prior protection of the labilizing amino group. This is commonly done by acetylation. Exhaustive methylation is one of a number of reactions of amines associated with the name of A. W. Hofmann. It involves a sequence of reactions terminating in the thermal decomposition of a quaternary ammonium hydroxide to yield a tertiary amine, usually trimethylamine, water, and an olefin. It has been widely used as a tool in the determination of the structures of complex compounds and, to a lesser extent, for the synthesis of olefins. The sequence of reactions is shown in (7). + AgOH 3CH3 l RCH2 CH2 NH2 −−− −−→ RCH2 CH2 N(CH3 )3 L− −−− −−→ + heat RCH2 CH2 N(CH3 )3 OH− −−− −−→ RCH
See ALKALOID.
CH + (CH3 )3 N + H2 O
(7)
Amino acid dating Preparation. Commercial preparation of aliphatic amines can be accomplished by direct alkylation of ammonia (Hofmann method, 1849) or by catalytic alkylation of amines with alcohols at elevated temperatures. Reduction of various nitrogen functions carrying the nitrogen in a higher state of oxidation also leads to amines. Such functions are nitro, oximino, nitroso, and cyano. For the preparation of pure primary amines, Gabriel’s synthesis and Hofmann’s hypohalite reaction are preferred methods. The Bucherer reaction is satisfactory for the preparation of polynuclear primary aromatic amines. Gabriel’s synthesis. This is a method for the synthesis of pure primary aliphatic amines by the hydrolysis of an N-alkyl phthalimide. The N-alkyl phthalimides are prepared by reaction of potassium phthalimide with an alkyl halide (preferably a bromide) [reaction (8)]. O
O
C
C NR
NK + RBr C
C
O
O
(8)
O C C
L-amino
OH
O Hofmann hypohalite reaction. This is another reaction associated with Hofmann, which furnishes a convenient method for the preparation of pure primary amines of either the aliphatic or aromatic series. In the overall sense, it involves conversion of an acid amide to an amine with loss of one carbon atom. The reaction proceeds through the stages shown in reactions (9)–(11). The amine arises by hydrolysis of RCONH2 + NaOH −→ RCONHX + NaOH
(9)
RCONHX + NaOH −→ O + NaX + H2 O
(10)
O + 2NaOH −→ RNH2 + Na2 CO3
(11)
R N
Determination of the relative or absolute age of materials or objects by measurement of the degree of racemization of the amino acids present. With the exception of glycine, the amino acids found in proteins can exist in two isomeric forms called D- and L-enantiomers. Although the enantiomers of an amino acid rotate plane-polarized light in equal but opposite directions, their other chemical and physical properties are identical. Amino acid handedness or homochirality is one of the most distinctive features of terrestrial life. It was discovered by L. Pasteur around 1850 that only L-amino acids are generally found in living organisms, but scientists still have not formulated a convincing reason to explain why life on Earth is based on only L-amino acids. See AMINO ACIDS; ENANTIOMER. Racemization. Under conditions of chemical equilibrium, equal amounts of both enantiomers are present (D/L = 1.0); this is called a racemic mixture. Living organisms maintain a state of disequilibrium through a system of enzymes that selectively utilize only the L-enantiomers. Once a protein has been synthesized and isolated from active metabolic processes, the L-amino acids are subject to a racemization reaction that converts them into a racemic mixture. The racemization reaction can be written as (1), where ki is the first-order rate constant for the
OH + RNH2
R
Amino acid dating
C
N
C
the isocyanate formed by migration of R from carbon to nitrogen in reaction (10) . Sodium hypobromite is the common laboratory reagent, but the cheaper calcium hypochlorite is used in commercial applications, such as the manufacture of anthranilic acid from phthalimide. See AMINO ACIDS; QUATERNARY Paul E. Fanta AMMONIUM SALTS. Bibliography. N. L. Allinger et al., Organic Chemistry, 2d ed., 1976; W. H. Brown and E. P. Rogers, General, Organic, and Biochemistry, 3d ed., 1987.
ki − D-amino acid acid − ki
(1)
interversion of the enantiomers. The rate of racemization is equal to 2ki. The kinetic equation of the racemization reaction is given by (2), where D/L is 1 + (D/L) 1 + (D/L) ln − ln = 2ki t (2) 1 − (D/L) 1 − (D/L) t=0 the ratio of the enantiomers of a certain amino acid at a particular time t. The t = 0 term is necessary to account for some racemization that occurs during sample preparation. Since racemization is a chemical process, the extent of racemization is dependent not only on the time that has elapsed since the L-amino acids were synthesized but also on the exposure temperature: the higher the temperature, the faster the rate of racemization. The rate of racemization is also different for most of the various amino acids. The half-life (for example, the time required to reach a D/L ratio of 0.33) at neutral pH for aspartic acid, an amino acid with one of the fastest racemization rates, is approximately 3500 years at 25◦C (77◦F) but is only 35 days at 100◦C (212◦F). Amino acids with the slowest racemization rates have half-lives roughly 10 times greater than those of aspartic acid. See HALF-LIFE. A variety of analytical procedures can be used to separate amino acid enantiomers; gas chromatography and high-performance liquid chromatography are the most widely used. Since these techniques have sensitivities in the parts-per-billion range, only a few hundred milligrams of sample material are
555
556
Amino acid dating normally required. Samples are first hydrolyzed in hydrochloric acid to break down the proteins into free amino acids, which are then isolated by cation-exchange chromatography. See CHROMATOGRAPHY. Since the late 1960s, the geochemical and biological significance of amino acid racemization has been extensively investigated. Geochemical uses of amino acid racemization include the dating of fossils or, in the case of known age specimens, the determination of their temperature history. Fossil types such as bones, teeth, and shells have been studied, and racemization has been found to be particularly useful for dating specimens that were difficult to date by other methods. Racemization has also been observed in the metabolically inert tissues of living mammals. Racemization can be studied in certain organisms and used to assess the biological age of a variety of mammalian species; in addition, it may be important in determining the biological lifetime of certain proteins. See RACEMIZATION. Dating fossils. Fossils have been found to contain both D- and L-amino acids, and the extent of racemization generally increases with geologic age. The amount of racemization in a fossil is determined by the value of ki and its age. The magnitude of ki in a fossil from a particular locality is primarily a function of a value known as the integrated exposure temperature, although protein diagenetic processes (for example, geochemical degradation and alteration) are also important. The rate constant ki can be estimated by using a calibration procedure, wherein the D/L ratio is measured in a fossil of known age from the study area. This D/L ratio and the calibration age of the sample are substituted into Eq. (2); an in-place ki value is thus calculated. This calibrated rate constant integrates variations in exposure temperature and other environmental parameters of the locality over the time period represented by the age of the calibration sample. Following calibration, the ki value can be used, with certain limitations, to date other samples from the surrounding region. The principal limitation is that the calibration sample should have a similar amino acid composition and content, and should be from roughly the same geologic period, as the sample being dated. The Olduvai Gorge region in the north-central Tanzanian Rift Valley offers an excellent opportunity to study the racemization reaction of amino acids in fossil bones and teeth over a time period extending back several million years. Since the geology of the region has been extensively studied as a result of the discovery of numerous hominid fossils and artifacts, the relative and absolute ages of the various stratigraphic units in the Olduvai region are well established. Thus, amino acid racemization can be investigated over several geologic periods. Aspartic acid racemization can be used to date only the Holocene and perhaps upper Pleistocene deposits in the Olduvai region. The ki values determined for aspartic acid at Olduvai suggest that bones and teeth older than about 60,000 years should contain only racemic aspartic acid, for example, D/L as-
partic acid = 1.0. However, fossil bones from the older stratigraphic units at Olduvai have been found to contain nonracemic aspartic acid. This indicates that these bones contain secondary aspartic acid, which probably is incorporated into fossil bone by percolating ground waters. Aspartic acid racemization of bones is of limited use in dating the important anthropological deposits in the Olduvai Gorge region because it racemizes rapidly and is prone to contamination problems. In contrast, the epimerization of L-isoleucine, which yields the nonprotein amino acid D-alloisoleucine, provides a much better chronological tool for dating the deposits in the Olduvai region. Isoleucine has two centers of asymmetry, and thus the reaction is termed epimerization rather than racemization; the equilibrium ratio of D-alloisoleucine to L-isoleucine is 1.3, and so Eq. (2) must be slightly modified to account for this. In fossil bones and teeth ranging in age from the upper Pleistocene to the Pliocene, the aile/ile ratio in tooth enamel steadily increases until an approximate equilibrium ratio of 1.3 is attained in the oldest deposits. This indicates that isoleucine epimerization in teeth is effectively a closed system with respect to the introduction of secondary isoleucine for a period of 3–4 million years. Isoleucine may be less susceptible to contamination than aspartic acid because the peptide bonds involving this amino acid are more stable with respect to hydrolysis than are those containing aspartic acid. By using the enamel sample from the base of Olduvai Bed I for calibration, isoleucine epimerization ages for the middle and lower Pleistocene deposits are obtained that are in good agreement with stratigraphic and radiometric age determinations. Some bones also yield reasonable ages. For example, the isoleucine epimerization age for a bone from the lower unit of the Masek beds is about 500,000 years, which is consistent with the age estimated from other evidence. Bones from older stratigraphic units, however, have been found to have aile/ile ratios less than those measured in tooth enamel. This implies that isoleucine in bones is more susceptible to contamination than is isoleucine in enamel. Isoleucine epimerization ages of bones from the older deposits should probably be viewed with caution and should be considered minimum age estimates. Contamination should consist primarily of L-amino acids; thus, incorporation of secondary amino acids into a fossil would lower the D/L ratio and yield an age estimate that is too young. The ki values for isoleucine epimerization in the upper Pleistocene deposits are about four times greater than those determined by using the lower Pleistocene Bed I sample. The faster rate of isoleucine epimerization in the younger deposits may be due to several factors, such as different exposure temperatures and the complex process of protein diagenesis. The upper Pleistocene ki value should be used only to date upper Pleistocene samples; it is not applicable to dating the older stratigraphic units.
Amino acid dating In general, the ages based on racemization and epimerization at Olduvai Gorge are consistent with those estimated from other evidence. The best results are obtained by using tooth enamel. On the basis of these results, the racemization and epimerization reactions in bones and teeth have been investigated at many localities throughout the world. Several hominid remains have been directly dated, and the technique has been used to evaluate dates derived from geologic, radiometric, and faunal evidence. Other racemization studies have utilized both terrestrial and marine fossil shells to estimate the age of various archeological sites as well as uplifted marine terrace deposits. In these cases, the limit for racemization dating is in the 100,000-year range, although at high latitudes where temperatures are cold the age range is extended nearly 1 million years. See FOSSIL. Racemization in living mammals. Enamel, dentine, and the proteins present in the nucleus of the eye lens are metabolically inert and are thus incubated at approximately 37◦C (98.6◦F) throughout a mammal’s life span. In mammals that live for decades, detectable racemization of aspartic acid has been found to take place in these tissues. Although the extent of racemization is quite small (for example, D/L ratios are in the range 0.02–0.10), it can be easily and precisely measured with available analytical techniques. Only aspartic acid racemization has been detected, which is expected since this amino acid has the fastest racemization rate. The racemization of amino acids thus can be used as a biochronological tool. This aging method is particularly useful for determining the ages of humans and certain other mammals, whose ages are not well documented. An interesting application of the racemization technique has been the determination of the ages of various marine mammals. For example, the ages of narwhals (Mondon monoceros), Arctic cetaceans that are characterized by the long (about 6.5 ft or 2 m) counterclockwise-spiraling tusk of the male, have been deduced by using the extent of aspartic acid racemization in their teeth. The method is particularly useful for determining the ages of female narwhals whose ages previously were largely unknown. The aspartic acid measurements indicate that both female and male narwhals can live more than 50 years. The method has also been used to estimate the ages of bowhead whales that live in the Arctic. During the recent hunting of these animals by Arctic peoples, some whales were found to have embedded in their skin stone harpoon tips, which were no longer used in the Arctic after around 1880. This finding suggested that bowhead whales might have unusually long life spans and, in order to help verify this surmise, a systematic study of the racemization in the eye lens nucleus was carried out. Extensive racemization was observed in the eye lens nucleus of some animals, and based on these results their ages were estimated to be in excess of 100 years. In the case of one whale, the calculated age was close
to 200 years. These findings suggest that bowhead whales may be some of the most aged mammals in the world. The extent of aspartic acid racemization in enamel and dentine can also be used to estimate the death age of humans from ancient burials. The main requirement in this application is that either the burial should be fairly recent or the burial environmental temperature should be relatively low, so that the extent of postmortem racemization is small compared with that which occurred during the lifetime of the individual. The racemization method has been used to calculate the age at death of mummified Alaskan Eskimos, of individuals in a medieval cemetery in Czechoslovakia, and of corpses in forensic cases. The racemization ages have greater accuracy (about ±5%) than those estimated by other means, and the method is especially useful in determining the ages of older individuals, whose age at death is often difficult to estimate from anatomical evidence. See ARCHEOLOGICAL CHEMISTRY; GEOCHRONOMETRY; RADIOCARBON DATING; ROCK AGE DETERMINATION. Other applications of amino acid racemization. Amino acid racemization provides a general indicator of the level of preservation of biomolecules in geological specimens. The extent of racemization of amino acids correlates with the level of overall degradation of labile molecules. One example deals with the preservation of ancient DNA. It has been found that only in specimens in which the extent of racemization is minimal (D/L < 0.1) are fragments of original DNA preserved. This has proved to be a valuable index for judging if DNA in fossils is original or derived from contamination. See DEOXYRIBONUCLEIC ACID (DNA). The reaction has implications with respect to protein turnover in mammalian tissues. It has been shown that there is an increasing level of waterinsoluble protein present in the eye lens nucleus, of long-lived mammals with increasing age. This is correlated with increasing racemization in the eye lens nucleus, and it has been suggested that structural changes induced by racemization result in protein insolubility. Amino acid racemization has also been found to take place during the processing of various food materials. This is caused by both heat exposure during cooking and the treatment of food components with alkali in order to make them more palatable. The racemization induced by these treatments may have deleterious nutritional effects because it has been found that proteins containing racemized amino acids are less digestible. The racemization of amino acids has implications in the search for life beyond Earth. It is assumed that life elsewhere would require homochiral amino acids, although whether it would be based on L- or D-amino acids is considered to be mainly a matter of a chance selection event that occurred either at the time of the origin of life or during early evolution. Detection of this amino acid homochirality signature would depend on whether life still existed on another world or, if it became extinct, whether racemization has erased the signature. In the case of
557
558
Amino acids Mars, using racemization rates determined in terrestrial environments, the survival time of an amino acid homochiral signature has been estimated for the conditions relevant to the surface of Mars throughout the planet’s history. Homochiral amino acids associated with a Martian biota that became extinct during the early history of the planet would still be preserved today in dry, cold (
vDS VBEn ≈ 0.65 V, Qp is off, that is, does not conduct (VEBp < 0), but Qn conducts and current flows from +VCC through Qn and the load resistance RL to ground: vo > 0. If vi < VEBp ≈ 0.65 V, the situation is reversed: Qn is off (VBEn < 0), and Qp conducts so that current flows from ground through RL and Qp to −VCC: vo < 0. Both transistors act as emitter followers, so that vo ≈ vi, but only one at a time conducts and no dc current flows from +VCC to −VCC. Apart from the small saturation voltage, VCE,sat ≈ 0.2 V, vo can swing from + VCC to −VCC with a gain of approximately unity and very little distortion. This behavior is depicted in the graph of vo versus vi, called the transfer curve (Fig. 11b), which also shows the socalled crossover distortion at the origin. This distortion arises because vi must at least equal the baseemitter voltages of the transistors before the emitter followers turn on and become active. The problem is avoided by arranging the circuit such that a dc voltage VBB equal to two forwardbiased pn-junction voltages, that is, ≈1.3 V always appears between the bases of the transistors (Fig. 11c). Two diodes D (in practice, two diode-connected transistors) are forward biased by a current source IB, which also supplies base currents to Qn and Qp. For ac signals, the forward-biased diodes are simply very small resistors. The operation is then as just described for the basic circuit (Fig. 11a), but the crossover distortion is avoided. The power conversion efficiency for this circuit, neglecting the small power dissipated in the diodes, can be shown to equal η = 0.785Vo–peak/VCC, that is, maximally 78.5%. By using the earlier results from the emitter follower, the gain is found to equal
Amplitude modulation approximately unity, and the input and output resistances are given by Eq. (39), where in and ip are the rin ≈ (β + 1)RL
rout ≈ ren ||rep =
VT in + i p
(39)
currents in the npn and pnp devices, respectively, and ren and rep are their emitter resistances as defined in Eq. (13). Thus, the input resistance is very large and the output resistance very small, as required for a good follower or buffer circuit. See POWER AMPLIFIER. Rolf Schaumann Bibliography. P. R. Gray and R. G. Meyer, Analysis and Design of Analog Integrated Circuits, 3d ed., 1993; R. Gregorian and G. C. Temes, Analog MOS Integrated Circuits for Signal Processing, 1986; A. S. Sedra and K. C. Smith, Microelectronic Circuits, 4th ed., 1997.
Amplitude (wave motion) The maximum magnitude (value without regard to sign) of the disturbance of a wave. The term “disturbance” refers to that property of a wave which perturbs or alters its surroundings. It may mean, for example, the displacement of mechanical waves, the pressure variations of a sound wave, or the electric or magnetic field of light waves. Sometimes in older texts the word “amplitude” is used for the disturbance itself; in that case, amplitude as meant there is called peak amplitude. This is no longer common usage. See LIGHT; SOUND. As an example, consider one-dimensional traveling waves in a linear, homogeneous medium. The wave disturbance y is a function of both a space coordinate x and time t. Frequently the disturbance may be expressed as y(x, t) = f(x ± υt), where υ denotes the wave velocity. The plus or minus sign indicates the direction in which the wave moves, and the shape of the wave dictates the functional form symbolized by f. Then, the amplitude of the disturbance at some point x0 is the maximum magnitude (that is, the maximum absolute value) achieved by f as time changes over the duration required for the wave to pass point x0. A special case of this is the one-dimensional, simple harmonic wave y(x, t) = A sin [k(x ± υt)], where k is a constant. The amplitude is A since the absolute maximum of the sine function is +1. The amplitude for such a wave is a constant. See HARMONIC MOTION; SINE WAVE. If the medium which a wave disturbs dissipates the wave by some nonlinear behavior or other means, then the amplitude will, in general, depend upon position. See WAVE MOTION. S. A. Williams
Amplitude modulation The process or result of the process whereby the amplitude of a carrier wave is changed in accordance with a modulating wave. This broad definition includes applications using sinusoidal carriers, pulse carriers, or any other form of carrier, the amplitude
factor of which changes in accordance with the modulating wave in any unique manner. See MODULATION. Amplitude modulation (AM) is also defined in a more restrictive sense to mean modulation in which the amplitude factor of a sine-wave carrier is linearly proportional to the modulating wave. AM radio broadcasting is a familiar example. At the radio transmitter the modulating wave is the audio-frequency program signal to be communicated; the modulated wave that is broadcast is a radio-frequency, amplitude-modulated sinusoid. See RADIO BROADCASTING. In AM the modulated wave is composed of the transmitted carrier, which conveys no information, plus the upper and lower sidebands, which (assuming the carrier frequency exceeds twice the top audio frequency) convey identical and therefore mutually redundant information. J. R. Carson in 1915 was the first to recognize that, under these conditions and assuming adequate knowledge of the carrier, either sideband alone would uniquely define the message. Apart from a scale factor, the spectrum of the upper sideband and lower sideband is the spectrum of the modulating wave displaced, respectively, without and with inversion by an amount equal to the carrier frequency. For example, suppose the audio-frequency signal is a single-frequency tone, such as 1000 Hz, and the carrier frequency is 1,000,000 Hz; then the lower and upper sidebands will be a pair of singlefrequency waves. The lower-sideband frequency will be 999,000 Hz, corresponding to the difference between the carrier and audio-signal frequencies; and the upper-sideband frequency will be 1,001,000 Hz, corresponding to the sum of the carrier and audiosignal frequencies. The amplitude of the signal appears in the amplitude of the sidebands. In practice, the modulating waveform will be more complex. A typical program signal might occupy a frequency band range of perhaps 100 Hz to 5000 Hz. An important characteristic of AM, as illustrated by Fig. 1, is that, apart from a scale factor and constant term, either the upper or lower envelope of the modulated wave is an exact replica of the modulating wave, provided two conditions are satisfied: first, that the carrier frequency exceeds twice the highest speech frequency to be transmitted; and second, that the carrier is not overmodulated. Single-sideband (SSB) modulation. This is a form of linear modulation that takes advantage of the fact that in an AM signal the carrier conveys no information and each of the two sidebands (that is, the upper sideband and the lower sideband) contains the same information. Hence it is not necessary to transmit both of them, and it is possible to achieve a 50% savings in bandwidth relative to normal AM or double-sideband (DSB) modulation. The price for this bandwidth efficiency is increased complexity. The generation of SSB is typically accomplished either by filtering out the unwanted sideband of a DSB signal or by using a phase-shift modulator. This latter system involves use of a Hilbert
605
606
Amplitude modulation
unmodulated carrier
3Ac 2
waveform of modulating wave
Ac Ac modulated carrier (a)
unmodulated carrier
Ac
waveform of modulating wave
Ac Ac modulated carrier
(b) unmodulated carrier
Ac 2
waveform of modulating wave
Ac Ac (c)
modulated carrier
Fig. 1. Amplitude modulation of a sine-wave carrier by a sine-wave signal, with (a) 50% overmodulation, (b) 100% modulation, and (c) 50% modulation. (After H. S. Black, Modulation Theory, Van Nostrand, 1953)
transformer to shift the phase of the frequency components of the message by −90◦. Specifically, if the message is denoted by m(t), the transmitted waveform s(t) is given by Eq. (1), where ω0 is the angular ˆ sin ω0 t s(t) = m(t) cos ω0 t ± m(t)
(1)
frequency of the carrier and m ˆ (t) is the Hilbert transform of m(t). If the plus sign is used in the above equation, the lower sideband is retained, whereas if the minus sign is used, the upper sideband is retained. Either a coherent demodulator or a carrier reinsertion receiver can be used to demodulate an SSB signal. In either case, there are certain limits on phase and frequency errors that can be tolerated. Also of concern in SSB transmission are linearity requirements, since the envelope of an SSB signal can undergo significant variations. SSB signaling is extensively used in voice transmission. Indeed, for applications such as telephone transmission over satellite links, one of the most common forms of modulation is the use of SSB subcarriers, which in turn frequency-modulate a carrier. The technique works as follows: Each voice message occupies roughly a 3-kHz bandwidth. The various voice messages are SSB-modulated onto individual subcarriers spaced 4 kHz apart. This composite baseband signal then frequency-modulates a carrier which is transmitted over the radio channel. SSB modulation is also used for voice transmission in radio relay systems and coaxial cable systems. Indeed, this type of transmission has become so well
accepted on a worldwide basis that standard frequency allocation plans have been set up by the Consultative Committee in International Telegraphy and Telephony (CCITT). The choice of whether or not to use SSB depends upon the particular conditions under which the system is operating. SSB requires both less power and less bandwidth than does AM, but requires a more complex transmitter and receiver. Also, SSB systems have been shown to be less sensitive to the effects of selective fading than are both AM and DSB. See SINGLE SIDEBAND. Vestigial-sideband (VSB) modulation. This is modulation whereby in effect the modulated wave to be transmitted is composed of one sideband plus a portion of the other adjoining the carrier. The carrier may or may not be transmitted. VSB is like SSB except in a restricted region around the carrier. The overall frequency response to the wanted sideband and to the vestigial sideband is so proportioned by networks that upon demodulation, preferably but not necessarily by a product demodulator, the original modulating wave will be recovered with adequate accuracy. By thus transmitting a linearly distorted copy of both sidebands in a restricted frequency region above and below the carrier, the original modulating wave can now be permitted to contain significant components at extremely low frequencies, even approaching zero in the limit. By this means, at the cost of a modest increase in bandwidth occupancy, network requirements are greatly simplified. Furthermore, the low-frequency limitation and the inherent delay associated with SSB are avoided. In standard television broadcasting in the continental United States the carrier is transmitted, envelope detection is used, and the vestigial sideband possesses a bandwidth one-sixth that of a full sideband. See TELEVISION. Comparison of modulation spectra. In order to better understand the differences between AM, DSB, SSB, and VSB, it is helpful to consider the resulting spectrum of each waveform when the signal given by Eq. (2) is used as the modulation, where A and B m(t) = A cos ω1 t + B cos ω2 t
(2)
are constant amplitudes. Figure 2a shows the singlesided spectrum of the modulation itself, where f1 = ω1/2π and f2 = ω2/2π; and Fig. 2b–e shows the spectra of the AM, DSB, SSB, and VSB waveforms, respectively, when modulated by m (t). In Fig. 2b, f0 is the carrier frequency and a is the modulation index. In Fig. 2e, represents the fraction of the lower sideband at frequency f0 − f1 which is retained, that is, that is, it represents a vestige of the lower sideband of the DSB waveform of Fig. 2c. Uses of AM in multiplexing. Multiplexing is the process of transmitting a number of independent messages over a common medium simultaneously. To multiplex, it is necessary to modulate. Two or more communicating channels sharing a common propagation path may be multiplexed by arranging them
Amplitude modulation A B
(a)
f2
f
f1
aB 2 (b)
aA 2
f0 − f2 f0 − f1
B 2
1
f0
aA 2
aB 2
f0 + f1 f0 + f2 A 2
A 2
f
f0
(c)
A 2
B 2 f
f0
(d)
A ( 1− ) 2
∋
A 2 (e)
f0
f
B 2
∋
along the frequency scale as in frequency division, by arranging them in time sequence as in time division, or by a process known as phase discrimination, in which there need be no separation of channels in either frequency or time. Frequency-division multiplexing. When communication channels are multiplexed by frequency division, a different frequency band is allocated to each channel. Single-sideband carrier telephone systems are a good example. At the sending end, the spectrum of each channel input is translated by SSB to a different frequency band. For example, speech signals occupying a band of 300–3300 Hz might be translated to occupy a band range of 12,300–15,300 Hz corresponding to the upper sideband of a sinusoidal carrier, the frequency of which is 12,000 Hz. Another message channel might be transmitted as the upper sideband of a different carrier, the frequency of which might be 16,000 Hz. At the receiving end, individual channels are separated by electric networks called filters, and each original message is recovered by demodulation. The modulated wave produced by SSB at the sending end becomes the modulating wave applied to the receiving demodulator. When communication channels are multiplexed by frequency division, all channels may be busy simultaneously and continuously, but each uses only its allocated fraction of the total available frequency range. Time-division multiplexing. When communication channels are multiplexed by time division, a number of messages is propagated over a common transmitting medium by allocating different time intervals in sequence for the transmission of each message. Figure 3 depicts a particularly simple example of a two-channel, time-division system. Ordinarily, the number of channels to be multiplexed would be considerably greater. Transmitting and receiving switches must be synchronized; time is of the essence in this system, and the problem is to switch at the right time. On the theoretical side there is a certain basic, fundamental question which must always be answered about any time-division system. The question is: At what rate must each communication channel be connected to its common transmitting path? Today it is known from the sampling principle that for successful communication each channel must be momentarily connected to the common path at a rate that is in excess of twice the highest message frequency conveyed by that channel. See PULSE MODULATION. Viewed broadly, amplitude modulation of pulse carriers generates the desired amplitude and phase relationships essential to time-division multiplexing. In addition, whereas each communication channel may use the entire available frequency band for transmitting its message, it may transmit only during its allocated fraction of the total time. Phase-discrimination multiplexing. This type of multiplexing, like SSB, saves bandwidth, may save signal power, and, like AM and VSB, has the important advantage of freely transmitting extremely low mod-
607
B 2 f
Fig. 2. Comparison of spectra of modulation waveforms. Signal amplitude is plotted as a function of frequency f. (a) Spectrum of modulation itself. (b) Normal amplitude modulation (AM). (c) Double-sideband (DSB). (d) Single-sideband (SSB). (e) Vestigial-sideband (VSB).
ulating frequencies. Furthermore, each communication channel may utilize all of the available frequency range all of the time. When n channels are multiplexed by phase discrimination, the modulating wave associated with each channel simultaneously amplitude-modulates n/2 carriers, with a different set of carrier phases provided for each channel. All sidebands are transmitted; n/2 carriers may or may not be transmitted. At the receiving end, with the aid of locally supplied carriers and an appropriate demodulation process, the n channels can be separated, assuming distortionless transmission, ideal modulators, and so on. Systems with odd numbers of channels can also be devised. Equality of sidebands and their exact phases account for the suppression of interchannel interference. Day’s system. This is a simple example of phasediscrimination multiplexing. Two sine-wave carriers of the same frequency but differing in phase by 90◦ are amplitude-modulated, each by a different message wave. The spectrum of each modulated sinusoid occupies the same frequency band. These modulated sinusoids are then added and propagated without distortion to a pair of demodulators. Quadrature carriers of correct frequency and phase are applied
608
Amplitude modulation
modulation
medium
modulation
(a) synchronous switching control
switching frequency
switching frequency
modulator demodulator
modulator demodulator
(b)
band-pass filter
F1 switching frequency
passes upper sideband synchronous switching control
F1 switching frequency
passes lower sideband
F2 (c)
modulator demodulator
F2 modulator demodulator
Fig. 3. Two-channel, time-division carrier system and two-channel, frequency-division carrier system. (a) General diagram of two-channel carrier system. (b) Amplitude modulation, time division. (c) Single-sideband modulation, frequency division. (After H. S. Black, Modulation Theory, Van Nostrand, 1953)
locally, one to each demodulator. Theoretically, a faithful copy of each message can be recovered. Within the continental United States, for purposes of saving bandwidth, Day’s system is used for multiplexing the two so-called color components associated with color television broadcasting. See MULTIPLEXING AND MULTIPLE ACCESS. Modulator and demodulator. Many methods of modulating and demodulating are possible, and many kinds of modulators and demodulators are available for each method. See MODULATOR. Fundamentally, since the sidebands of an amplitude-modulated sinusoid are generated by a multiplication of wave components which produces frequency components corresponding to the products, it is natural to envisage a product modulator having an output proportional to the product of two inputs: modulating wave and carrier. An ideal product modulator suppresses both modulating wave and carrier, transmitting only upper and lower sidebands. See AMPLITUDE MODULATOR. At the receiving end, the original message may be
recovered from either sideband or from both sidebands by a repetition of the original modulating process using a product modulator, commonly referred to as a product demodulator, followed by a lowpass filter. Perfect recovery requires a locally applied demodulator carrier of the correct frequency and phase. For example, a reduced carrier system creates its correct demodulator carrier supply by transmitting only enough carrier to control the frequency and phase of a strong, locally generated carrier at the receiver. See AMPLITUDE-MODULATION DETECTOR. SSB systems commonly generate their modulator and demodulator carriers locally with independent oscillators. For the high-quality reproduction of music, the permissible frequency difference between modulator and demodulator carriers associated with a particular channel is limited to about 1–2 Hz. For monaural telephony, frequency differences approaching 10 Hz are permissible. Unbalanced square-law demodulators, rectifiertype demodulators, and envelope detectors are often used to demodulate AM. However, even though
Amplitude-modulation detector overmodulation (Fig. 1) is avoided, significant distortion may be introduced. In general, the amount of distortion introduced in this manner will depend upon the kind of demodulator or detector used, the amount of noise and distortion introduced prior to reception, and the percentage modulation. Harold S. Black; Laurence B. Milstein Bibliography. L. W. Couch, Digital and Analog Communication Systems, 5th ed., 1996; M. Schwartz, Information, Transmission, Modulation and Noise, 4th ed., 1990; H. Stark, F. B. Tuteur, and J. B. Anderson, Modern Electrical Communications, 2d ed., 1988; F. G. Stremlers, Introduction to Communication Systems, 3d ed., 1990; H. Taub, Principles of Communication Systems, 2d ed., 1986; R. E. Ziemer and W. H. Tranter, Principles of Communications, 4th ed., 1994.
Amplitude-modulation detector A device for recovering information from an amplitude-modulated (AM) electrical signal. Such a signal is received, usually at radio frequency, with information impressed in one of several forms. The carrier signal may be modulated by the information signal as double-sideband (DSB) suppressed-carrier (DSSC or DSBSC), double-sideband transmittedcarrier (DSTC or DSBTC), single-sideband suppressed-carrier (SSBSC or SSB), vestigial sideband (VSB), or quadrature-amplitude modulated (QAM). The field of amplitude-modulation detector requirements splits by application, complexity, and cost into the two categories of synchronous and asynchronous detection. Analog implementation of nonlinear asynchronous detection, which is typically carried out with a diode circuit, is favored for consumer applications, AM-broadcast radio receivers, minimum-cost products, and less critical performance requirements. Synchronous detectors, in which the received signal is multiplied by a replica of the carrier signal, are implemented directly according to their mathematics and block diagrams, and the same general detector satisfies the detection requirements of all SSB, DSB, and VSB signals. Although synchronous detectors may operate in the analog domain by using integrated circuits, more commonly digital circuits are used because of cost, performance, reliability, and power advantages. A quadrature-amplitude-modulation detector is simply a two-channel version of a synchronous detector. Synchronous detection. The idea behind synchronous detection is simple. There are two conceptual approaches: to reverse the modulation process (which is rather difficult), or to remodulate the signal from the passband (at or near the transmitter’s carrier frequency) to the baseband (centered at dc or zero frequency). The remodulation approach is routine. Unfortunately, a nearly exact replica of the transmitter’s carrier signal is needed at the receiver in order to synchronously demodulate a transmitted signal. Synchronous means that the carrier-signal reference in the receiver has the same frequency and
phase as the carrier signal in the transmitter. There are three means available to obtain a carrier-signal reference. First, the carrier signal may actually be available via a second channel. There are no difficulties with this method because a perfect replica is in hand. Second, the carrier signal may be transmitted with the modulated signal. It then must be recovered by a circuit known as a phase-lock loop with potential phase and frequency errors. Third, the carrier signal may be synthesized by a local oscillator at the receiver, with great potential for errors. Unless otherwise stated, it will be assumed that a perfect replica of the carrier signal is available. See OSCILLATOR; PHASE-LOCKED LOOPS. SSB, VSB, and DSB signals. Demodulation of an SSBmodulated (upper or lower), VSB-modulated, or DSBmodulated signal is as simple as multiplying it by the transmitter carrier-signal replica and low-pass filtering to remove unwanted high-frequency terms (sometimes called images) from the product in order to obtain the base-band signal (Fig. 1a). Sometimes the multiplication is carried out in a single step, sometimes in two steps called conversion stages. If the type of transmitted-signal modulation is transmitted carrier (TC), the consequent dc component must be removed by a filter. If the information signal itself has a dc component, an amplitude-modulated TC method cannot be used to transmit the data because of the difficulty in separating the dc component of the signal from the carrier which is shifted to dc by the demodulation process. See ELECTRIC FILTER.
four-quadrant multiplier modulated input signal
demodulation reference= carrier replica (a)
filter
reference source
baseband signal
baseband signal+ high-frequency products
four-quadrant multiplier
QAM input signal
I-channel reference
filter
I-channel output
filter
Q- channel output
reference source −j
−90° phase shifter Q-channel reference
four-quadrant mutiplier (b) Fig. 1. Synchronous (coherent) detectors for (a) doublesideband (DSB), single-sideband (SSB), and vestigialsideband (VSB) signals, and (b) for quadrature-amplitudemodulated (QAM) signals.
609
610
Amplitude-modulation detector If the ith information component in a modulated signal is Ai cos (wit) and the modulating carrier is cos (wct), then the ith modulated element is given by Eq. (1). The ideal local oscillator signal is y(t) = x(t) = Ai cos (wi t) cos (wc t)
(1)
2 cos (wct), so the demodulated signal is given by Eq. (2). The last term is removed by the filter, so that z(t) = x(t)y(t) = 2Ai cos (wi t) cos2 (wc t) = 2(Ai /2) cos (wi t)[1 + cos (2wc t)] = Ai cos (wi t) + [Ai cos (wi t) cos (2wc t)] (2) the information signal is completely recovered. See TRIGONOMETRY. The center frequency of the signal is often shifted down in two conversion stages. First an analog conversion by an oscillator translates the frequency of the signal to a value convenient for the analog-todigital converter to operate; then the four-quadrant digital multiplier serves as a second converter and completes the frequency translation. The increasing performance and decreasing costs of analog-todigital converters favor detectors in which analog conversion is eliminated and the analog-to-digital converter operates directly on the higher-frequency signal. The importance of the accuracy of the reference in synchronous demodulation cannot be overemphasized. An average frequency error in the demodulating reference is intolerable because it will translate the modulated signals to the wrong frequencies. The cosine of the phase angle in the demodulating reference is the demodulator gain. A phase error can therefore seriously reduce the gain. A small frequency error integrates with time to become a
phase error, causing objectionable periodic amplitude modulation of the received signal. However, it is not necessary for the reference waveform to be a pure sinusoid, as long as it does have the correct periodicity and phase. An appeal to Fourier theory shows that this reference waveform can be represented by a sum of sinusoids harmonically related to the fundamental frequency of the reference. Multiplying the signal to be detected by this nonsinusoidal reference causes so-called images to appear after the multiplication process. These images, which exist at multiples of the fundamental frequency, can be removed by suitably filtering the multiplier output signal. See FOURIER SERIES AND TRANSFORMS; NONSINUSOIDAL WAVEFORM. QAM signals. A synchronous detector of quadrature-amplitude-modulated signals consists of two synchronous demodulators operating in parallel (Fig. 1b). The two demodulating carrier-reference signals are separated in phase by precisely 90◦. The accuracy of the phase angle of the carrier-reference signal limits the degree of separation (after demodulation) that can be maintained between the two QAM signals that have been transmitted on top of one another. If the ith transmitted signal element is given by Eq. (3), and the I and Q demodulation carrier references are given by Eqs. (4) and (5), then the filtered x(t) = Ai cos (wi t) cos (wc t) +Bi cos (wi t) sin (wc t) (3) yI (t) = 2 cos (wc t + ) = 2[cos (wc t) cos ( ) − sin (wc t) sin ( )] (4) yQ (t) = 2 sin (wc t + ) = 2[sin (wc t) cos ( ) + cos (wc t) sin ( )] (5) (that is, the double-frequency products have been
rectifier
source
filter
Q
P R1
R2
C3
C1
R4 R3
AGC voltage
half-wave
full-wave
Fig. 2. Diode detector circuit. P = input point of filter; Q = output point of filter.
C2
output to audio amplifier and filter
Amplitude-modulation detector Diode detector. A diode detector circuit in a radio receiver has three stages (Fig. 2). The first stage is the signal source, which consists of a pair of tuned circuits that represent the last intermediate-frequency (i.f.) transformer which couples the signal energy out of the i.f.-amplifier stage into the detector. The second stage is a diode rectifier, which may be either full wave or half wave. Finally, the signal is passed through the third stage, a filter, to smooth high-frequency noise artifacts and remove the average value. See DIODE; INTERMEDIATE-FREQUENCY AMPLIFIER; RADIO RECEIVER; RECTIFIER. The waveform shaping at the input point of the filter (Fig. 2) is determined by a capacitor C1 in parallel with an equivalent resistance of R1, the latter in series with a parallel combination of resistors, R2 and R4. The filter also has a capacitor C3 between R2 and R4, and a parallel combination of a resistor R3 and a capacitor C2 in series with R2. The reactances of both capacitors C2 and C3 are quite small at the information frequency, so the capacitors can be viewed as short circuits. Meanwhile, the C3−R4
removed) demodulated outputs are given by Eqs. (6) and (7). Separation is complete for = 0. zI (t) = Ai cos (wi t) cos ( ) − Bi cos (wi t) sin ( ) (6) zQ (t) = Bi cos (wi t) cos ( ) + Ai cos (wi t) sin ( ) (7) Asynchronous detection. Asynchronous detection applies to DSTC signals whose modulation index is less than 1, and is quite simple. First the received signal is full-wave rectified, then the result is lowpass filtered to eliminate the carrier frequency and its products, and finally the average value is removed (by dc blocking, that is, ac coupling). This result is identical to that which is obtained by synchronous detection with a reference that has been amplitude distorted to a square wave. The cheaper but less efficient half-wave rectifier can also be used, in which case the demodulation process is called envelope detection.
amplitude
2
0
−2
time
(a) 2
amplitude
amplitude
2
0
0
time
time (c)
(b)
output signal
information signal
−1 (d)
1
amplitude
amplitude
1
output signal
information signal
−1
time
time (e)
Fig. 3. Detection of a 50%-modulated double-sideband transmitted-carrier (DSTC) signal. An artificially low frequency ratio of 10:1 is used to illustrate raggedness of the output. (a) Received signal and modulating information-signal envelope. (b) Full-wave rectified received signal. (c) Half-wave rectified signal. (d) Information signal and output signal after full-wave rectifier and filter. (e) Information signal and output signal after half-wave rectifier and filter.
611
612
Amplitude-modulation radio combination serves as a dc-blocking circuit to eliminate the constant or dc component at the output point of the filter. In order to bias the output point properly for the next amplifier stage, R4 is replaced by a biasing resistor network in a real filter; R4 is the single-resistor equivalent. See ELECTRICAL IMPEDANCE. The strength of the signal arriving at the detector (Fig. 3a) is proportional to the mean value of the signal at the input point of the filter and is sensed as the automatic-gain-control (AGC) voltage. Changes in this voltage level are used to adjust the amplification before the detector so that the signal strength at the detector input can remain relatively constant, although the signal strength at the receiver’s antenna may fluctuate. Since capacitor C2 shunts all signal energy and any surviving carrier energy to ground, the AGC voltage is roughly the average value at the input point of the filter scaled by R3/(R1 + R2 + R3) because R1 is much smaller than R4. See AUTOMATIC GAIN CONTROL (AGC). Waveforms. The received signal (Fig. 3a) is delivered from the i.f. stage to the diode rectifier. If the filter were simply a resistor, then the output from the diode rectifier would simply be the full-wave rectification (Fig. 3b) or half-wave rectification (Fig. 3c) of the received signal. (The presence of the capacitors in the filter provides the waveform shaping.) In both cases the modulating (information) signal shapes the envelopes of the rectified waveforms. In an amplitude-modulation radio, the output from the filter after full-wave rectification (Fig. 3d) would be the audio signal (plus some noise) that is delivered to the audio amplifier. Additional filtering can be provided as necessary to reduce the noise to an acceptable level. This amplified and filtered signal is finally delivered to an output device, such as a loudspeaker. The ragged waveform of the filter output contrasts with the smooth waveform of the information signal (Fig. 3d). The raggedness vanishes as the ratio of the i.f. frequency to the information frequency increases. While a synthetic example with a low ratio can be used to clearly show the effects within the demodulator (Fig. 3d), the amplitude of this raggedness noise decreases in almost direct proportion to the increase in the frequency ratios. The raggedness of the output signal from the filter after half-wave rectification (Fig. 3e) is much greater than that of the full-wave-rectifier case. The actual worst-case ratio for standard-broadcast amplitude-modulation radio is 46.5:1. In this case the raggedness on the filtered outputs is reduced to a fuzz that can be seen in graphs of the waveforms but is well outside the frequency range of audio circuits, loudspeakers, and human hearing. See AMPLITUDE MODULATION; AMPLITUDE MODULATOR; MODULATION; WAVEFORM. Stanley A. White Bibliography. R. C. Dorf (ed.), The Electrical Engineering Handbook, 2d ed., 1997; M. S. Roden, Analog and Digital Communication Systems, 3d ed., 1991; M. Valkenburg (ed.), Reference Data for Engineers: Radio, Electronics, Computer, and Communications, 8th ed., 1996.
Amplitude-modulation radio Radio communication employing amplitude modulation of a radio-frequency carrier wave as the means of conveying the desired intelligence. In amplitude modulation the amplitude of the carrier wave is made to vary in response to the fluctuations of a sound wave, television image, or other information to be conveyed. See AMPLITUDE MODULATION; RADIO. Amplitude modulation (AM), the oldest and simplest form of modulation, is widely used for radio services. The most familiar of these is broadcasting; others include radiotelephony and radiotelegraphy, television picture transmission, and navigational aids. The essentials of these radio systems are discussed in this article. Low frequency (long wave). European and Asian countries use frequencies in the range 150–255 kHz for some broadcast services. An advantage of these frequencies is stable and relatively low-attenuation wave propagation. When not limited by atmospheric noise, large areas may be served by one station. In the United States these frequencies are reserved for navigational systems and so are not available for broadcasting. See RADIO SPECTRUM ALLOCATION; RADIOWAVE PROPAGATION. Low-frequency (lf) broadcast antennas are omnidirective and radiate vertically polarized waves. Unless special means are used to reduce antenna selectivity, the highest program frequencies transmitted are substantially below 10,000 Hz. Medium frequency. The frequencies in the range from 535 to 1705 kHz are reserved for AM (standard) broadcasting. In the Western Hemisphere this band is divided into channels at 10-kHz intervals, certain channels being known as clear, regional, and local, according to the licensed coverage and class of service. The local channels are occupied by stations, usually of 250-W output, servicing smaller localities. Many stations occupy a channel, but they are situated far enough apart to permit interference-free coverage in the local area. Fewer stations of higher power, but with greater distances between them, share the regional channels. A few clear channels are occupied by high-power stations (50,000-W maximum output in the United States). These stations may have exclusive use of a channel, or may share it with another distant station. See RADIO BROADCASTING. Interference between co-channel regional stations and clear-channel stations is minimized by use of directive antennas, which suppress radiation toward other stations and direct it to main populated areas from a properly located station. European medium-frequency (mf) broadcasting channels are assigned at 9-kHz intervals rather than the 10-kHz intervals used in the Western Hemisphere. This reduced spacing provides more channels within the mf band. The technique of directive antennas, which also provides more channels within a band by minimizing interference between stations, has not been used extensively in Europe. Vertically polarized radiation is used at medium and low frequencies propagated over the Earth’s
Amplitude-modulation radio surface. There is also propagation of high-angle radiation via reflection from the ionosphere, a phenomenon that predominates at night but is relatively absent during daylight. This sky-wave propagation accounts for the familiar long-distance reception at night. At distances where the downcoming sky waves overlap the ground wave, fading and distortion of the signal occurs. Receivers in the groundwave zone get stable signals day and night. In the overlap zone, daylight reception may be stable but night reception disturbed by fading. In the sky-wave zone at night, assuming no interference from other stations, satisfactory service may be received over long distances. Reception in this zone depends on atmospheric noise, which varies seasonally, and the state of the ionosphere, which varies greatly from night to night depending upon astronomical conditions that affect the upper atmosphere. Individual AM broadcast stations transmit program frequencies ranging from 30 to 10,000 Hz with excellent fidelity. To obtain suitable tuning selectivity between channels, AM broadcast receivers may reproduce less extensive program frequencies, according to make, cost, type, and condition of the receiver. High frequency (shortwave). Small bands of frequencies between 3000 and 7500 kHz are used in tropical areas of high atmospheric noise for regional broadcasting. This takes advantage of the lower atmospheric noise at these frequencies and permits service under conditions where medium frequencies have only severely limited coverage. Wave propagation day and night is by sky wave. Groundwave coverage from such stations is usually negligible, since high-angle horizontally polarized radiation is used. Short-distance coverage by this mode is by ionospheric reflection of waves radiated almost vertically. Long-distance international broadcasting uses high-power transmitters and directive (beam) antennas operating in the bands 5950–6200 kHz, 9500– 9775 kHz, 11,700–11,975 kHz, 15,100–15,450 kHz, 17,700–17,900 kHz, 21,450–21,750 kHz, and 25,600–26,100 kHz. These bands are allocated throughout the world for this purpose, and the band used for any path depends upon ionospheric conditions which vary with direction, distance, hour, season, and the phase of the 11-year sunspot cycle. Typically, waves are propagated in low-angle beams by multiple reflections between ionosphere and Earth to cover transoceanic distances, and signals are often distorted during transmission. These bands are so crowded that a signal is seldom received without interference for any prolonged period. Reception from particular stations can be improved by the use of special directive receiving antennas. Propagation partially through the ionosphere also causes rotation of the transmission’s plane of polarization due to the Faraday effect. See FARADAY EFFECT. The technical performance of high-frequency (hf) broadcast transmission systems is usually to the same standards employed for mf broadcasting, although propagation and interference conditions seldom make this evident to a listener.
AM telephony and telegraphy. The first radiotelephony was by means of amplitude modulation, and its use has continued with increasing importance. Radiotelephony refers to two-way voice communication. Amplitude modulation and a modified form called single-sideband are used almost exclusively for radiotelephony on frequencies below 30 MHz. Above 30 MHz, frequency or phase modulation is used almost exclusively, a notable exception being 118–132 MHz, where amplitude modulation is used for all two-way very high frequency (vhf) radiotelephony in aviation operations. The least expensive method known for communicating by telephony over distances longer than a few tens of miles is by using the high frequencies of 3–30 MHz. Furthermore, since radio is the only way to communicate with ships and aircraft, hf AM radiotelephony has remained essential to these operations, except for short distances that can be covered from land stations using the very high frequencies. Therefore, hf radiotelephony has become established for a great variety of pioneering and exploration enterprises and private, public, and government services needing telephone communication, fixed or mobile, over substantial distances where there are no other ways to telephone. Because AM techniques are simple and inexpensive, this form of modulation has predominated. The economic development of distant hinterland areas depends greatly on hf AM radio telephony to the outside world, although satellite transmission and reception has generally replaced hf AM radio telephony for all purposes except broadcasting. Widespread use has led to serious crowding of the hf band. All governments are suspending the use of the hf band wherever it is technically and economically feasible to employ frequencies above 30 MHz using either direct transmission or radio repeater stations. The trend to single-sideband modulation also alleviates the pressure of congestion in the hf band. Many of the radiotelephone systems in use operate two ways on one frequency in simplex fashion, that is, all stations on the frequency are normally in a receiving status. Transmission, by press-to-talk (manual) or voice-operated carrier (automatic) switching from reception to transmission, occurs only while the sender is talking. Many stations can thus occupy one frequency provided the volume of traffic by each of the stations is small. Two-way telephony must be strictly sequential. This system is not adapted for connection to a normal two-wire telephone. Full duplex radiotelephony, for interconnection with wire telephone systems, is essential for most public correspondence. This requires two frequencies, each available full time in one direction. Even so, typical fading of signals during propagation requires that voice-operated antiregeneration devices be used to maintain circuit stability. Talking must be sequential between speakers as there can be no interrupting, but the system will interconnect with conventional business or home telephones. Amplitude-modulated telegraphy consists of
613
614
Amplitude-modulation radio interrupting a carrier wave in accordance with the Morse dot-dash code or codes used for the printing telegraph. Much of the radiotelegraphic traffic of the world uses AM telegraphy, although there has been extensive conversion to frequency-shift (frequency-modulation) telegraphy since 1944, the latter being better adapted to automatic teleprinting operations. Radiotelegraph operations have been refined, speeded, and mechanized, but, under adverse noise and fading conditions, AM manual telegraphy between experienced operators is still more reliable. Most aviation and marine radiotelegraphy uses AM manual methods. See TELEGRAPHY; TELEPHONE SERVICE. Single-sideband (SSB) hf telephony. This is a modified form of amplitude modulation in which only one of the modulation sidebands is transmitted. In some systems the carrier is transmitted at a low level to act as a pilot frequency for the regeneration of a replacement carrier at the receiver. Where other means are available for this purpose, the carrier may be completely suppressed. Where intercommunication between AM and SSB systems is desired, the full carrier may be transmitted with one sideband. Since 1954 the use of SSB has expanded rapidly in replacing common AM telephony for military and many nonpublic radio services. In time, SSB will gradually displace AM radiotelephony to reduce serious interference due to overcrowding of the radio spectrum. SSB transmission also is less affected by selective fading in propagation. Multiplexing, both multiple-voice channels or teleprinter channels included with voice, is applied to SSB transmission. Teleprinting and data transmission by frequency multiplex using SSB radiotelephone equipment is increasing and displacing older methods of radiotelegraphy for fixed point-to-point government and common-carrier services. See SINGLE SIDEBAND. Aviation and marine navigation aids. Amplitudemodulated radio has a dominant role in guidance and position location, especially in aviation, which is almost wholly under radio guidance. Radio is used in traffic control from point to point, in holding a position in a traffic pattern, and in approach and landing at airports. Distance measuring from a known point, radio position markers, runway localizers, and glidepath directors all use AM in some form, if only for coded identification of a facility. Marine operations are not so dependent on radio facilities as are those of aviation, but almost every oceangoing ship makes use of direction finding, at least to determine its bearing from a radio station of known location. Special marine coastal beacon stations emit identified signals solely for direction-finding purposes. Certain navigational systems (Decca, loran) for aviation are also used by ships for continuous position indication and guidance. See ELECTRONIC NAVIGATION SYSTEMS. Television broadcasting. Amplitude modulation is used for the broadcasting of the picture (video) portion of television. In England, France, and a few other places amplitude modulation is also used for the
sound channel associated with the television picture, but frequency modulation is more commonly used for sound. Countries of the Western Hemisphere, Japan, Philippines, Thailand, and Iran broadcast television video in an emission band of 4.25 MHz; the English video bandwidth is 3 MHz; the French system, 10 MHz. The rest of continental Europe (except Russia) use a bandwidth of 5.25 MHz. The carrier frequencies employed are between 40 and 216 MHz, and 470 to 890 MHz. A channel allocation includes the spectrum needed for both sound and picture. Japan and Australia also use 88–108 MHz for television broadcasting. Reception of such transmissions is often impaired by multipath (the same signal arriving out-of-phase at the receiver from different propagation paths due to reflections by large objects and terrain), causing effects such as ghosting. The English and French systems employ positive video modulation, in which white corresponds to higher amplitudes of the modulation envelope and black corresponds to lower amplitudes. All other established video broadcasting uses negative modulation, working in the opposite sense. Synchronizing pulses in the negative system are at maximum carrier amplitude. The dynamic range from black to white in the picture varies from 75 to 25% of maximum amplitude. The upper-frequency portion of one video sideband is suppressed by filters so that its remaining vestige, together with the other complete sideband, is transmitted. This is called vestigial sideband transmission and avoids unnecessary spectrum usage. See TELEVISION TRANSMITTER. Edmund A. Laport; Michael C. Rau Stereophonic broadcasting. A number of systems for stereophonic AM radio broadcasting have existed since the late 1970s. These systems allow AM radio stations to transmit two channels of information in the same spectrum space where only one could exist previously. The result is similar to that of stereophonic FM, audio cassette, and other binaural entertainment media. To transmit in AM stereo, a broadcast station employs a device known as an exciter to adapt its existing transmitter. The exciter has left- and rightchannel audio inputs as well as summed audio- and radio-frequency outputs. Stereophonic audio program material is connected to the audio inputs, while the outputs attach to the AM transmitter. The summed audio (left channel plus right channel) information is used to amplitude-modulate the transmitter, to assure that compatibility with existing monophonic receivers is maintained. The radio-frequency exciter output, however, contains the stereophonic information to be transmitted. This output is connected in place of the transmitter’s crystal oscillator, and although the signal present is at the exact carrier frequency of the radio station, it is also phasemodulated by the difference (left channel minus right channel) audio information. A subaudible and low-level identification tone is usually added to this difference audio information. This allows AM radio
Amplitude modulator receivers to identify stereophonic transmissions and activate the appropriate decoding circuitry, as well as a stereo indicator lamp commonly found on the radio front panel. See PHASE MODULATION; RADIO RECEIVER. Stereophonic AM broadcasting methods all have their roots in a system known as quadrature multiplexing. This system allows the transmission of two channels of information on a single carrier frequency, but true quadrature transmissions are inherently incompatible with the majority of radios available to consumers. The differences between various systems employed to transmit stereophonic AM broadcasts all relate to the methods used to overcome this incompatibility. See AMPLITUDE MODULATION; RADIO; STEREOPHONIC RADIO TRANSMISSION. Stanley Salek Bibliography. National Association of Broadcasters, NAB Guide to Advanced Television Service, 2d ed., 1991; J. Whitaker (ed.), NAB Engineering Handbook, 9th ed., National Association of Broadcasters, 1999.
Amplitude modulator A device for moving the frequency of an information signal, which is generally at baseband (such as an audio or instrumentation signal), to a higher frequency, by varying the amplitude of a mediating (carrier) signal. The motivation to modulate may be to shift the signal of interest from a frequency band (for example, the baseband, corresponding to zero frequency or dc) where electrical disturbances exist to another frequency band where the information signal will be subject to less electrical interference; to isolate the signal from shifts in the dc value, due to bias shifts with temperature or time of the characteristics of amplifiers, or other electronic circuits; or to prepare the information signal for transmission. Familiar applications of amplitude modulators are standard-broadcast or amplitude-modulation (AM) radio; data modems (modem = modulator + demodulator); and remote sensing, where the information signal detected by a remote sensor is modulated by the remote signal-conditioning circuitry for transmission to an electrically quiet data-processing location. See AMPLITUDE-MODULATION RADIO; MODEM; REMOTE SENSING. The primary divisions among amplitude modulators are double sideband (DSB); single sideband (SSB); vestigial sideband (VSB); and quadrature amplitude (QAM), where two DSB signals share the same frequency and time slots simultaneously. Each of these schemes can be additionally tagged as suppressed carrier (SC) or transmitted carrier (TC). To amplitude-modulate an information signal is (in its simplest form) to multiply it by a second signal, known as the carrier signal (because it then carries the information). A real (as opposed to complex) information signal at baseband, or one whose spectrum is centered about zero frequency, has an amplitude spectrum which is symmetric (an even func-
f
0 (a)
f (b) upper sideband
upper sideband
lower sidebands
−fc
0
fc
−fc
0
fc
−fc
0
fc
−fc
0
fc
f
(c)
f
(d)
f
(e)
f
(f) Fig. 1. Baseband spectra and spectra of outputs from suppressed-carrier modulators. Amplitudes and phases are shown as functions of frequency f. (a) Baseband amplitude spectrum. (b) Baseband phase spectrum. (c) Amplitude spectrum of output from double-sideband (DSB) modulator. (d) Amplitude spectrum of output from upper singlesideband (SSB) modulator. (e) Amplitude spectrum of output from lower single-sideband (SSB) modulator. (f) Amplitude spectrum of output from vestigial-sideband (VSB) modulator.
tion; Fig. 1a) and a phase spectrum which is asymmetric (an odd function; Fig. 1b) in frequency about zero. Because of this symmetry, the information in the upper or positive sideband (positive frequencies) replicates the information in the lower or negative sideband (negative frequencies). Balanced modulation or linear multiplication (that is, four-quadrant multiplication where each of the two inputs is free to take on positive or negative values without disturbing the validity of the product) of this information signal by a sinusoidal carrier of single frequency, fc, moves the information signal from baseband and replicates it about the carrier frequencies fc and −fc (Fig. 1c). See FOURIER SERIES AND TRANSFORMS. Suppressed-carrier modulators. Linear multiplication (described above) produces double-sideband suppressed-carrier (DSSC or DSBSC) modulation of the carrier by the information. One of the redundant sidebands of information may be eliminated by filtering (which can be difficult and costly to do adequately) or by phase cancellation (which is usually a much more reasonable process) to produce the more
615
616
Amplitude modulator efficient single-sideband suppressed-carrier (SSBSC or SSB) modulation (Fig. 1d, e). The process of only partially removing the redundant sideband (usually by deliberately imperfect filtering) and leaving only a vestige of it is called vestigial-sideband (VSB) modulation (Fig. 1f). See ELECTRIC FILTER. The DSB and SSB generation mechanisms may be illustrated by considering one sinusoid, Ai cos (wit), of a weighted sum of sinusoids that may be assumed to constitute the information signal, u(t). The ith information signal component can be multiplied by the carrier signal, cos (wct), to produce the DSSC product, x(t), as in Eq. (1). This expresx(t) = Ai cos (wi t) cos (wc t) A1 {cos [(wc − wi )t] + cos [(wc + wi )t]} (1) = 2 sion clearly shows the generation of the lower-sideband-frequency signal component at wc−wi and the upper-sideband-frequency signal component at wc + wi. See TRIGONOMETRY. A pair of phase shifters may now be introduced to operate on both the carrier and information signals to produce sin (wct) and Ai sin (wit), respectively. The product of this new signal pair is y(t), given by Eq. (2), which exhibits sideband-frequency signal y(t) = Ai sin (wi t) sin (wc t) =
A1 {cos [(wc − wi )t] − cos [(wc + wi )t]} (2) 2
components at the same frequencies as x(t). However, the sum of x(t) and y(t) is given by Eq. (3), x(t) + y(t) = Ai cos [(wc − wi )t]
(3)
which is the lower-sideband-frequency signal only, and the difference of x(t) and y(t) is given by Eq. (4), x(t) − y(t) = Ai cos [(wc + wi )t]
(4)
which is the upper-sideband-frequency signal only. Implementation of this phase-shifting method to generate either an upper- or a lower-SSB signal re-
four-quadrant multiplier
carrier generator
information signal wideband −90° phase shifter
narrowband −90° phase shifter
four-quadrant multiplier Fig. 2. Single-sideband (SSB) amplitude modulator.
+ + −
∑
SSB - modulated signal output
quires a pair of linear modulators (Fig. 2), one to generate x(t), the other to generate y(t). The phase shifting of the information signal is generally performed by a wideband Hilbert transformer, a filter whose gain is −j for positive frequencies and j for negative frequencies, where j is the square root of −1. The narrowband phase shifter for the carrier signal can be a simple second-order filter. However, such a filter may not be necessary because the carrier-frequency generator often provides both sine and cosine outputs simultaneously. Because the required transmission bandwidth is 50% less than that of DSB, SSB is an exemplary modulation system for bandwidth preservation. See SINGLE SIDEBAND. Since the redundant sideband has been reduced (but not eliminated) in the generation of a VSB signal, its transmission efficiency is less than that of a SSB signal but greater than that of a DSB signal. The VSB signal can be generated in two steps: manufacture of a DSB signal, followed by the partial filtering out of a redundant sideband. Carrier insertion. Receivers demand a carrier-signal reference to properly demodulate (detect) the transmitted signal. This reference may be provided in one of three ways: it may be transmitted with the modulated signal, for example, in double-sideband transmitted-carrier (DSTC or DSBTC) modulation, which is the method used for standard-broadcast AM radio; transmitted on a separate channel (often the case for instrumentation systems); or generated at the receiver. Receivers for double-sideband signals can be very inexpensive if the carrier is transmitted with the information-bearing sidebands, as in the case of DSTC signals. In concept, the carrier may be inserted by simply adding it to the DSSC signal. For maximum transmission efficiency in generating the DSTC signal, the minimum amount of carrier should be transmitted. (Transmission efficiency may be defined as the ratio of total sideband power to the sum of sideband-plus-carrier power.) However, in order to avoid distortion in simple low-cost receivers, the amplitude of the inserted carrier must be at least equal to the peak amplitude, A, of the information signal. The ratio of the actual amplitude of the information signal to A is known as the modulation index, which should never exceed unity (Fig. 3). See AMPLITUDEMODULATION DETECTOR; DISTORTION (ELECTRONIC CIRCUITS). The transmitter can also be simplified by adding to the information signal a dc value equal to the peak magnitude of the information signal. The value of this modified information signal is now always positive. No longer is a four-quadrant multiplier required to modulate; rather, a far simpler two-quadrant multiplier generates the DSTC signal directly. Mechanization. DSSC modulation, no matter how it is disguised, is just ordinary (linear) multiplication, or a reasonable approximation to that multiplication. Furthermore, as discussed above, a DSTC signal can be modeled as a DSSC signal with the carrier added. Any amplitude modulator is therefore simply some sort of embodiment of a linear multiplier
Amplitude modulator
3
information signal waveform amplitude
amplitude
3
0
−3
0
−3 time
(a) 3
1
amplitude
amplitude
time
(b)
0
−3
0
−1 time
(c)
time (d)
Fig. 3. Waveforms of a signal after various types of double-sideband modulation. Information signal envelopes are shown superimposed on the modulated waveforms. Double-sideband transmitted-carrier (DSTC) modulation with modulation indices of (a) 50%, (b) 100%, and (c) 200% are compared with (d) double-sideband suppressed-carrier (DSSC) modulation.
with or without a means to add the carrier. High-level DSTC modulation is usually done by varying the power-supply voltage to the final highpower amplifier in the radio-frequency transmitter by summing the information signal with the dc output voltage of the power supply. Low-level DSTC modulation is carried out by integrated circuits that approximate the multiplications, followed by a linear amplifier. See AMPLIFIER; ELECTRONIC POWER SUPPLY. Because of the falling costs and improving performance of digital devices, signals are now generally represented by digital data, that is, signal values which are sampled in time and encoded as numbers to represent their sampled values. Digital signalprocessing (DSP) functions are carried out by some sequence of addition (or subtraction), multiplication, and delays. Efficient mechanization of each of these functions has been the subject of considerable effort. Digitally, waveform generation and linear multiplication are highly optimized processes. Advances in integrated-circuit fabrication techniques have so lowered the cost of digital circuits for modulation that the overwhelming majority of amplitude modulators in use are now digital. Custom applicationspecific integrated circuits (ASICs), general-purpose programmable DSP devices, and customizable arrays such as programmable logic arrays (PLAs) and fieldprogrammable arrays (FPAs) are all in widespread use as amplitude modulators and demodulators. Digital modems as data transmission equipment dominate the production of amplitude modulators. See INTEGRATED CIRCUITS.
Quadrature amplitude modulation (QAM). It was observed above that the DSSC signal could be reduced to an SSB signal because of the redundant information in the second sideband. This redundancy can be exploited in a different way by quadrature modulation. In addition to the first information signal, u(t), quadrature amplitude modulation acts on a second information signal, v(t). The purpose of the process is to amplitude-modulate the two signals to the same frequency, wc, with the goal of recovering each signal intact and uncontaminated by the other. When a first DSSC signal, u(t) cos (wct), is demodulated with a first reference signal, cos (wct), the product signal given by Eq. (5) is obtained. After low-pass u(t)[cos (wc t)]2 = 0.5[1 + cos (2wc t)]u(t)
(5)
filtering to eliminate the double-frequency term, cos (2wct), and scaling by 2 to offset the 0.5 gain, the first information signal, u(t), is recovered. If the first received signal, u(t) cos (wct), had been demodulated with a second reference signal, sin (wct), the product signal given by Eq. (6) would have been obu(t) sin (wc t) cos (wc t) = 0.5 sin (2wc t)u(t)
(6)
tained, and low-pass filtering this signal to eliminate the double-frequency term, sin (2wct), would leave nothing. Similarly, when a second transmitted and received DSSC signal, v(t) sin (wct), is demodulated with the second reference signal, sin (wct), the product signal
617
618
Amylase given by Eq. (7) is obtained. After low-pass filterv(t)[sin (wc t)]2 = 0.5[1 − cos (2wc t)]v(t)
(7)
ing to eliminate the first double-frequency term, cos (2wct), and scaling by 2 to offset the 0.5 gain, the second information signal, v(t), is obtained. If the second received signal, v(t) sin (wct), had been demodulated with the first reference signal, cos (wct), the product signal given by Eq. (8) would have been v(t) sin (wc t) cos (wc t) = 0.5 sin (2wc t)v(t)
is found almost exclusively in higher plants. See CARBOHYDRATE; ENZYME; GLYCOGEN.
(8)
obtained, and low-pass filtering this signal to eliminate the double-frequency term, sin (2wct), would again leave nothing. The two carrier signals, sin (wct) and cos (wct), are 90◦ apart; they are also orthogonal, or are in quadrature. That is, the average value of their product is zero. An obvious deduction is that, if a third signal is transmitted that is the sum of the first and second transmitted signals of the discussion above, that is, u(t) cos (wct) + v(t) sin (wct), and the received signal is applied to two parallel channels (named by convention I and Q), then it is possible to demodulate in the first (I) channel with the first reference signal, cos (wct), and in the second (Q) channel with the second reference signal, sin (wct), thereby recovering the first and second information signals, u(t) and v(t), respectively, from the I and Q channels. The modulation method in the discussion above could be DSBTC if u(t) and v(t) were replaced by u (t) = K1 + u(t) and v (t) = K2 + v(t), where K1 and K2 are equal to or greater than the maximum amplitudes of u(t) and v(t), respectively, in order to keep the modulation indices positive. In the above discussion, the first and second transmitted signals are both amplitude-modulated signals, but the modulating signals [cos (wct) and sin (wct)] are in quadrature. The pair of information signals, u(t) and v(t), are therefore quadrature-amplitudemodulated (QAM) in the formation of the third transmitted signal. See AMPLITUDE MODULATION; MODULATION; MODULATOR. Stanley A. White Bibliography. R. C. Dorf (ed.), The Electrical Engineering Handbook, 2d ed., 1997; M. S. Roden, Analog and Digital Communication Systems, 3d ed., 1991; M. Valkenburg (ed.), Reference Data for Engineers: Radio, Electronics, Computer, and Communications, 8th ed., 1996.
Amylase An enzyme which breaks down (hydrolyzes) starch, the reserve carbohydrate in plants, and glycogen, the reserve carbohydrate in animals, into reducing fermentable sugars, mainly maltose, and reducing nonfermentable or slowly fermentable dextrins. Amylases are classified as saccharifying (β-amylase) and as dextrinizing (α-amylases). The α- and β-amylases are specific for the α- and β-glucosidic bonds which connect the monosaccharide units into large aggregates, the polysaccharides. The α-amylases are found in all types of organs and tissues, whereas β-amylase
Animals. In animals the highest concentrations of amylase are found in the saliva and in the pancreas. Salivary amylase is also known as ptyalin and is found in humans, the ape, pig, guinea pig, squirrel, mouse, and rat. Pig pancreas is rich in amylase, whereas cattle, sheep, and dog pancreases have lower concentrations. Starch is one of the most important constituents of human food. The food prepared by the mouth for swallowing (the bolus) is converted by the gastric juices into chyme (a semiliquid paste). Chyme is passed through the pylorus into the duodenum, where intestinal digestion occurs. Part of this digestion is caused by pancreatic amylase, which, like ptyalin, hydrolyzes starch to maltose. A maltase also found in pancreatic juice hydrolyzes maltose to glucose. Glucose is picked up by the bloodstream for use in the tissues for respiration and for conversion to glycogen in the liver for storage. Plants. Starch is broken down during the germination of seeds (rich in starch) by associated plant enzymes into sugars. These constitute the chief energy source in the early development of the plant. β-Amylase occurs abundantly in seeds and cereals such as malt. It also is found in yeasts, molds, and bacteria. Industry. Amylase is also used in industry. It is used (1) in brewing and fermentation industries for the conversion of starch to fermentable sugars, (2) in the textile industry for designing textiles, (3) in the laundry industry in a mixture with protease and lipase to launder clothes, (4) in the paper industry for sizing, and (5) in the food industry for preparation of sweet syrups, to increase diastase content of flour, for modification of food for infants, and for the removal of starch in jelly production. The amylase enzyme used in industry comes from many sources: fungi, malt, bacteria, and the pancreas gland of cattle. See MALT BEVERAGE. Daniel N. Lapedes
Amyloidosis A disorder characterized by the accumulation of an unusual extracellular fibrous protein (amyloid) in the connective tissue of the body. The deposition of amyloid may be widespread, involving major organs and leading to serious clinical consequences, or it may be very limited with little effect on health. The term amyloid originated with R. Virchow (1854), who mistakenly believed that this substance was akin to starch. Amyloid protein has now been found to have the following unique electron microscopic, x-ray diffraction, and biochemical characteristics. It appears as a homogeneous, eosinophilic extracellular substance when stained with hematoxylin and eosin. When it is stained with Congo red, green birefringence patterns are visible in the polarizing microscope. With the electron microscope, it appears in the form of fine, nonbranching fibrils 7–10 nanometers in diameter, and x-ray diffraction
Amyloidosis studies show a cross beta pattern. Amino acid sequence studies demonstrate the presence of unusual proteins (AA, AL transthyretin, beta-2 microglobin, beta/A4 protein, and others). Amyloidosis has been classified clinically as: (1) primary amyloidosis, with no evidence for preexisting or coexisting disease; (2) amyloidosis associated with multiple myeloma; (3) secondary amyloidosis, associated with chronic infections (such as osteomyelitis, tuberculosis, leprosy), chronic inflammatory disease (such as rheumatoid arthritis, ankylosing spondylitis, regional enteritis), or neoplasms (such as medullary carcinoma of the thyroid); (4) heredofamilial amyloidosis, associated with familial Mediterranean fever and a variety of heritable neuropathic, renal, cardiovascular, and other syndromes; (5) local amyloidosis, with local, often tumorlike, deposits in isolated organs without evidence of systemic involvement; (6) amyloidosis associated with aging in the heart or brain; (7) amyloid associated with hormonal disorders; and (8) amyloid of chronic hemodialysis. Immunochemical findings. There is an immunochemical classification based on definition of the various types of proteins in amyloidosis: (1) amyloid associated with tissue deposition of protein AA, a tissue protein occurring in secondary amyloidosis and the amyloid of familial Mediterranean fever; and (2) amyloid associated with tissue deposition of protein AL, a protein consisting primarily of the aminoterminal variable segment of the light chain of homogeneous immunoglobulins, or in a few instances the whole light chain, which occurs in primary amyloid and amyloid associated with multiple myeloma. Transthyretin, the major protein of autosomal dominant inherited neuropathic amyloid, is also found in the amyloid lesions in the aged heart. Other chemical types also exist; for example, beta protein is found in the amyloid of Alzheimer’s disease, IAPP (amylin) in that of adult onset (type 2) diabetes, and beta-2 microglobulin in the amyloid of chronic hemodialysis. Antisera to protein AA react with an antigenically related component, serum AA (SAA), that is found in normal human serum, but has a higher molecular weight than AA and is regarded as the precursor of AA. Serum AA has several distinct characteristics: (1) It behaves as an acute phase reactant and is elevated in concentration in infection and inflammation. (2) It migrates as an apolipoprotein in the serum. (3) It is produced in the liver. In the heredofamilial amyloidosis, a variety of single point mutations have been found as the cause of the disorder; that is, the prototype is transthyretin with the substitution of methionine for valine at position 30. A second component of amyloid, P-component (plasma component or pentagonal unit), with different ultrastructure, x-ray diffraction pattern, and chemical characteristics, has been isolated from amyloid and shown to be identical with a circulating human alpha globulin present in serum in minute amounts. It has certain similarities to C-reactive protein (CRP, an acute-phase protein in humans). How-
ever, it does not behave as an acute-phase protein in humans, and has biophysical and immunologic distinctions from CRP. Clinical features. The different types of amyloid overlap considerably with regard to organ distribution. Since any organ may be involved, the clinical manifestations of amyloidosis are extremely varied. Infiltration of and about peripheral nerves, skin, tongue, joints, and heart is most frequent in the primary forms, while in secondary amyloidosis the liver, spleen, and kidney are the major sites of deposits. Carpal tunnel syndrome (median neuropathy), resulting from local deposition of amyloid, is a frequent finding in certain types of primary amyloidosis and in the amyloidosis associated with myeloma and some heredofamilial forms. Some individuals with primary amyloidosis develop a coagulation disorder in which there is a deficiency of factor X (a procoagulant in plasma). This condition has been traced to the binding of these proteins to amyloid fibrils. Amyloid in the heredofamilial (involving multiple siblings of successive generations) forms is responsible for peripheral neuropathy in certain families in Portugal, Japan, and Sweden; upper limb neuropathy in certain families of German or Swiss ancestry in the United States; cranial nerve neuropathy in a group of Finnish ancestry; and kidney disease and cardiovascular disease in others. All of the above conditions are inherited as autosomal dominant traits except for the amyloid associated with familial Mediterranean fever, which is inherited in an autosomal recessive manner. Clinical features in the diagnosis of primary amyloidosis include unexplained proteinuria, peripheral neuropathy, enlargement of the tongue, enlargement of the heart, intestinal malabsorption, bilateral or familial carpal tunnel syndrome, or low blood pressure when standing upright. Symptoms frequently exist for several years before correct diagnosis. When primary amyloidosis is suspected, evidence of an underlying disease is investigated to rule out possible inflammatory disorders or malignant tumors. In cases of chronic infectious or inflammatory diseases, secondary amyloidosis is suspected if there is unexplained proteinuria, or enlargement of the liver or spleen. The diagnosis is established by means of microscopic examination of appropriate biopsy specimens. Subcutanous abdominal fat or tissue from the rectum have proved to be particularly useful, and in selected cases, skin, gingival, kidney, and liver biopsy specimens are valuable. Since a biopsy specimen may contain only a small quantity of amyloid, it is important to look for the characteristic green birefringence on polarization microscopy after Congo red staining. Treatment. There is no specific treatment for amyloidosis, but supportive treatment is very useful. Some individuals with secondary renal amyloidosis have improved after the cure of a predisposing disease, such as active tuberculosis or chronic osteomyelitis. Kidney transplants have been carried out successfully in a few cases of renal failure. Colchicine
619
620
Anaerobic infection has been shown to relieve the acute attacks of familial Mediterranean fever and may decrease the formation of amyloid. Combinations of melphalan, prednisone, and colchicine are used in primary (AL) amyloidosis, and liver transplantation has been introduced in heredofamilial amyloid to remove the site of synthesis of the mutant protein. Alan S. Cohen Bibliography. A. S. Cohen, Amyloidosis, Bull. Rheumatic Dis., 40:1–12, 1991; A. S. Cohen, Amyloidosis, N. Engl. J. Med., 277:522–530, 1967; J. R. Harris (ed.), Electron Microscopy of Proteins, vol. 3, 1982; J. B. Natvig et al., Amyloid and Amyloidosis, 1990.
Anaerobic infection An infection caused by anaerobic bacteria (organisms that are intolerant of oxygen). Most such infections are mixed, involving more than one anaerobe and often aerobic or facultative bacteria as well. Anaerobes are prevalent throughout the body as indigenous flora, and virtually all anaerobic infections arise endogenously, the principal exception
being Clostridium difficile colitis. Factors predisposing to anaerobic infection include those disrupting mucosal or other surfaces (trauma, surgery, and malignancy or other disease), those lowering redox potential (impaired blood supply, tissue necrosis, and growth of nonanaerobic bacteria), drugs inactive against anaerobes (such as aminoglycosides), and virulence factors produced by the anaerobes (toxins, capsules, and collagenase, hyaluronidase, and other enzymes). Anaerobic gram-negative bacilli (Bacteroides, Prevotella, Porphyromonas, Fusobacterium) and anaerobic gram-positive cocci (Peptostreptococcus) are the most common anaerobic pathogens. Clostridium (spore formers) may cause serious infection. The prime pathogen among gram-positive nonsporulating anaerobic bacilli is Actinomyces. Infections commonly involving anaerobes and the incidence of anaerobes in such infections are noted in the table. In terms of frequency of occurrence, the oral and dental pleuropulmonary, intraabdominal, obstetric-gynecologic, and skin and soft tissue infections are most important. Foul odor is definitive evidence of involvement of anaerobes in an
Infections commonly involving anaerobic bacteria
Infections
Incidence, %
Proportion of cultures positive for anaerobes yielding only anaerobes
Bacteremia Central nervous system Brain abscess Extradural or subdural empyema Head and neck Chronic sinusitis Chronic otitis media Neck space infections Wound infection following head and neck surgery Dental, oral, facial Orofacial, of dental origin Bite wounds Thoracic Aspiration pneumonia Lung abscess Empyema (nonsurgical) Abdominal Intraabdominal infection (general) Appendicitis with peritonitis Liver abscess Other intraabdominal infection (postsurgery) Biliary tract Obstetric-gynecologic Miscellaneous types Pelvic abscess Vulvovaginal abscess Vaginal cuff abscess Septic abortion, sepsis Pelvic inflammatory disease Soft tissue and miscellaneous Nonclostridial crepitant cellulitis Pilonidal sinus infection Diabetic foot ulcers Soft-tissue abscesses Cutaneous abscesses Decubitus ulcers with bacteremia Osteomyelitis
20
4/ 5
89 10
1/ 2 to 2/ 3
52 33 to 56 100 95
4/ 5∗ 0 to 1/ 10 3/ 4 0
94 47
4/ 10 1/ 34
62 to 100 85 to 93 62 to 76
1/ 3 to 1/ 2† 1/ 2 to 3/ 4 1/ 3 to 1/ 2
81 to 94 96 52 93 41 to 45
1/ 10 to 1/ 3 1/ 100 1/ 3 1/ 6 2/ 117 to 0
71 to 100 88 75 98 63 to 67 25 to 48
1/ 3 1/ 2 1/ 4 1/ 30
∗
75 73+ 95 60 62 63 40
1/ 14 to 1/ 7 1/ 12 1/ 20 1/ 4 1/ 5 1/ 10
23/ 28 cultures (82%) yielding heavy growth of one or more organisms had only anaerobes present. † Aspiration pneumonia occurring in the community, rather than in the hospital, involves anaerobes to the exclusion of aerobic or facultative forms two-thirds of the time. SOURCE: S. M. Finegold, Antimicrobial therapy of anaerobic infections: A status report, Hosp. Pract., 14:17—81, October 1979.
Analgesic infection, but absence of such odor does not exclude the possibility. Other important clues to anaerobic infection, though not specific, are location near a mucosal surface; bite infection; tissue necrosis; gas formation; prior aminoglycoside, quinolone, or trimethoprim/sulfamethoxazole therapy; “sulfur granules” (actinomycosis); certain distinctive clinical presentations (such as gas gangrene); septic thrombophlebitis; unique morphology of certain anaerobic organisms on smear; and failure to grow organisms aerobically from specimens. To document anaerobic infection properly, specimens for culture must be obtained so as to exclude normal flora and must be transported under anaerobic conditions. See GANGRENE; GAS. Therapy includes surgery (debridement, drainage) and antimicrobial agents. Bacteroides fragilis is the most resistant anaerobe; only chloramphenicol, carbapenems, beta lactam/beta lactamase inhibitor combinations, and metronidazole are essentially active against all strains. Other useful drugs include clindamycin, cefoxitin, ceftizoxime, and cefotetan. See ANTIBIOTIC; INFECTION. Sydney M. Finegold Bibliography. S. M. Finegold and W. L. George (eds.), Anaerobic Infections in Humans, 1989; P. N. Levett (ed.), Anaerobic Microbiology: A Practical Approach, 1992; T. Willis and L. Phillips, Anaerobic Infections: Clinical and Laboratory Practice, 1991.
Analcime A mineral with a framework structure in which all the aluminosilicate tetrahedral vertices are linked, thus allying it to the feldspars, feldspathoids, and zeolites. The structure is cubic, space group Ia3d, α = 1.371 mm. Its formula is Na(H2O)[AlSi2O6]; in this sense it is a tectosilicate. The analcime structure type includes several other mineral species. These include high-temperature leucite, β-K[AlSi2O6]; pollucite, Cs(H2O)-[AlSi2O6]; and wairakite, Ca(H2O)[AlSi2O6]2. Temperaturedependent order-disorder relationships between AlO4 and SiO4 tetrahedra lead to modifications of lower structural symmetry, birefringent optical prop-
31 mm (a)
(b)
Analcime, (a) Specimen from Keweenaw County, Michigan (American Museum of Natural History). (b) Trapezohedral crystal typical of analcime (after C. Klein, Manual of Mineralogy, 21st ed., John Wiley and Sons, 1993).
erties, and complex composite twinning about a single crystal nucleus. Crystals are most often trapezohedra (see illus.); rarely the mineral is massive granular. Hardness is 5–5.5 on Mohs scale; specific gravity is 2.27. Analcime most frequently occurs as a lowtemperature mineral in vesicular cavities in basalts, where it is associated with zeolites (particularly natolite), datolite, prehnite, and calcite. Small grains are frequent constituents of sedimentary rocks and muds in oceanic basins associated with volcanic sources. Noted localities for large crystals include the Bay of Fundy region, Nova Scotia; the Watchung basalts of New Jersey; and the Keweenaw Peninsula in northern Michigan. Analcime also occurs in certain hydrothermal sulfide deposits as a fissure mineral. See FELDSPAR; FELDSPATHOID; SILICATE MINERPaul B. Moore ALS; ZEOLITE.
Analgesic Any of a group of drugs of diverse chemical structure and physiological effects which are commonly used for the relief of pain. To qualify as an analgesic a drug must selectively reduce or abolish pain without causing impairment of consciousness, mental confusion, incoordination or paralysis, or other derangements of the nervous system. Narcotic alkaloids. The oldest and best-known analgesics are opium, a drug obtained by extracting the juice of the poppy seed, and its most active alkaloid, morphine. Morphine and related drugs reduce or block the activation of pain neurons in the gray matter of the spinal cord, and at receptor sites in the brainstem and thalamus. The mechanisms by which this inhibition is effected are complex. Morphine does not act directly on the pain-mediating neurons. Instead, it involves an elaborate pain-modulating system in the posterior horns of the spinal cord, the dorsal-ventral reticular formation of the medulla, the periaqueductal gray matter, and the hypothalamus and thalamus. Separate neurons in these regions have opioid receptors on their surface, and the strength of this affinity between drug and neuron is proportional to the analgesic potency of the opiate (or opioid). All parts of this neuronal system acting locally or by way of descending connections inhibit pain transmission cells. In these same regions there are also neurons that elaborate a naturally occurring (endogenous) group of opioid peptides whose function is to modulate or suppress pain. They are called enkephalins, and the most potent of them is betaendorphin. Since they are abundant not only in these sites but also in the gut, in sympathetic neurons, and in chromafin cells of the adrenal glands, they are believed to account for the natural inhibition of pain by strong emotion (stress), and expectancy of pain relief by nonanalgesic drugs (placebo effect). See ENDORPHINS; MORPHINE ALKALOIDS; OPIATES. In addition to their use as analgesic drugs, opiates have other biological effects such as sedation,
621
622
Analog computer pupillary constriction, suppression of cough reflex, respiratory depression, reduction of intestinal motility, impairment of segmental flexor reflexes, and decrease in body temperature. This functional diversity is attributed to the activation of other inhibitory systems of neurons. See SEDATIVE. While morphine is the most powerful medical analgesic substance, there are many other naturally occurring alkaloids derived from opium. The best known of these is codeine. Common to all opiates is the attribute that if they are taken for weeks or months the recipient will need larger doses to obtain the same analgesic and sedative effects. This response is called tolerance. If the drug is stopped, disagreeable withdrawal or abstinence effects are experienced within hours to days. There is severe pain, sweating, salivation, hyperventilation, restlessness, and confusion. These abstinence symptoms, which are marks of habituation, pressure the addicted person to take extreme measures to obtain the narcotic in order to avoid the symptoms. See ADDICTIVE DISORDERS; ALKALOID; NARCOTIC. Synthetics. Because of the strong addictive properties of morphine and related compounds, chemists have synthesized other drugs of similar chemical structure, in the hope of securing analgesia without addiction. This effort has been only partially successful. Methadone, a drug that has been given to addicts as a substitute for morphine, is an effective analgesic when taken orally and is less addictive than morphine. Meperidine (Demerol) is a strong synthetic analgesic but definitely addictive. Other synthetic analgesics are oxycodone (Percodan), levorphanol (levodromoran), propoxyphene (Darvon), and pentazocine (Talwin). The last two of this series cause little or no addiction but, unfortunately, are not strong analgesics. Another synthetic drug, Naloxone, blocks the analgesic effect of all opiate agonists and precipitates withdrawal symptoms in addicted individuals. Salicylates. Another class of analgesic drugs, which are nonnarcotic (nonaddictive), are the salicylates, the most familiar being acetysalicylic acid (aspirin), and salicylatelike drugs such as phenylbutazone (Butazolidine), indomethacin (Indocin), acetaminophen, and phenacetin. These drugs are most effective in relieving skeletal pain due to inflammation (such as arthritis). Their analgesic properties, which are not nearly as strong as those of morphine and the synthetic opioids, are due to their action on both the peripheral and central nervous system. Peripherally they block the synthesis of prostaglandins in the inflamed tissues and thereby prevent the sensitization of pain receptors. Centrally they act in some obscure way on the hypothalamus. These drugs also have many other effects, such as reducing fever (antipyrexia) and preventing platelet agglutination. They are the most commonly used of all analgesic medications and are often combined with caffeine or a barbiturate sedative under a variety of trade names and sold for the relief of headache, backache, and so forth. See ASPIRIN; EICOSANOIDS; NERVOUS SYSTEM (VERTEBRATE); PAIN. Raymond D. Adams
Analog computer A computer or computational device in which the problem variables are represented as continuous, varying physical quantities. An analog computer implements a model of the system being studied. The physical form of the analog may be functionally similar to that of the system, but more often the analogy is based solely upon the mathematical equivalence of the interdependence of the computer variables and the variables in the physical system. See SIMULATION. Types. An analog computer is classified either in accordance with its use (general- or specificpurpose) or based on its construction (hydraulic, mechanical, or electronic). General-purpose implies programmability and adaptability to different applications or the ability to solve many kinds of problems. Most electronic analog computers were general-purpose systems, either real-time analog computers in which the results were obtained without any significant time-scale changes, or high-speed repetitive operation computers. The latter utilized time compression typically in the range 1000:1 to 10,000:1, completing a computation or simulation in a few milliseconds and enabling the display of results as flicker-free images. Special-purpose analog computers contain fixed programs permitting few or no adjustments. They are generally built into physical systems where they serve typically as a subsystem simulator, function controller or results analyzer. For example, the pneumatic computer shown in Fig. 1 used air bellows and flapper nozzles to generate accurate multiplication, division, squaring, or square-root functions of input signals, encoding data as air pressures. Use. Large electronic analog computer systems with many hundreds of operational amplifiers were widely used from the early 1960s to the mid-1980s. They solved extremely complex and extensive sets of differential equations (mathematical models) such as six-degree-of freedom space flights, exothermal chemical reaction kinetics, control systems for food processing plants, and the human immunosuppressive system. During the 1970s, they were often paired with a digital computer in hybrid systems
Fig. 1. Special-purpose pneumatic analog computer, model Foxboro 556. (Foxboro Co.)
Analog computer
Fig. 2. Analog-hybrid laboratory equiped with a HYSHARE 600 system. This system can include up to six analog computer consoles and one or more digital processors for multiuser, multitask applications. (Electronic Associates, Inc.)
(Fig. 2). The digital computer then primarily controlled the setting up and running of the analog computers, collecting and checking results and often providing logic and initial values to allow extensive searching through vast multidimensional problem domains. In some instances, hybrid systems also entailed automatic patching (interconnecting) of analog computing units. Smaller analog computers, with fewer than 100 operational amplifiers, were extensively used for education and training purposes. Models of a theoretical nature, expressed by sets of differential equations, or practical engineering simulations based on empirical data and heuristic models often lent themselves to study by analog computers. The ease by which the user could interact with the modeled system was a major feature, since learning is greatly enhanced when the user safely can explore, experiment with, and even misuse the modeled system. Digital equivalents. Since the 1970s, digital computer programs have been developed which essentially duplicate the functionality of the analog computer. Modern simulation languages, such as ACSL, GASP, GPSS, SLAM, and Simscript, have replaced electronic analog computers. They provide nearly the same highly interactive and parallel solution capabilities of electronic analog computers, but without the technical shortcomings of electronics: accuracy inherently limited to 0.01%, effective bandwidths of 1 MHz, and cumbersome and time-consuming programming. Simulation languages also avoid the large purchase investments and the continual maintenance dependencies of complex electronic systems.
Digital multiprocessor analog system. Another type of analog computer is the digital multiprocessor analog system (Fig. 3), in which the relatively slow speeds of sequential digital increment calculations have been radically boosted through parallel processing. In this type of analog computer it is possible to retain the programming convenience and data storage of the digital computer while approximating the speed, interaction potential, and parallel computations of the traditional electronic analogs. The digital multiprocessor analog computer typically utilizes several specially designed high-speed processors for the numerical integration functions, the data (or variable) memory distributions, the arithmetic functions, and the decision (logic and control) functions. All variables remain as fixed or floatingpoint digital data, accessible at all times for computational and operational needs. The digital multiprocessor analog computer achieves an overall problem-solving efficiency comparable to the very best continuous electronic analog computers, at a substantially lower price. An example of such a computer, the model AD10 (Fig. 3a), can solve large, complex, and multivariate problems at very high speeds and with the advantages of alldigital hardware. Its computation system (Fig. 3b) is based on five parallel processors working in a special computer architecture designed for high-speed operation. The various elements are interconnected by means of a data bus (MULTIBUS). A highly interleaved data memory of up to 106 words serves the data and program storage functions. The five processors working in parallel are: the control processor (COP), which controls all operations; the arithmetic processor (ARP), which runs the numerical
623
624
Analog computer
DEP decision processor program memory 1024 32-bit words
MAP memory address processor program memory 1024 48-bit words
(a)
17 banks of 4096 64 -bit words
ARP arithmetic processor program memory 1024 80-bit words
address/control multibus (AM) data multibus (DM) status multibus (SM)
multibus
IOCC input /output channel controller 128 external devices*
interleaved data memory
COP control processor
HIC host interface controller
program memory 1024 32-bit words
*analog - to - digital converters, digital - to - analog converters, digital - to - digital interfaces
future processors and input /output
NIP numerical integration processor program memory 1024 80-bit words
RIC remote interface control PDP II host processor
(b) Fig. 3. Multiprocessor analog computer. (a) Exterior. (b) Organization. (Applied Dynamics International)
calculations; the decision processor (DEP), which executes the logic parts of the program; the memory address processor (MAP), which makes sure that all data are fetched and stored efficiently; and the numerical integration processor (NIP), which carries out the integration functions that are crucially important in an analog computer. Description. The typical electronic generalpurpose analog computer consists of a console containing a collection of operational amplifiers; computing elements, such as summing networks, integrator networks, attenuators, multipliers, and function generators; logic and interface units; control circuits; power supplies; a patch bay; and various meters and display devices. The patch bay is arranged to bring input and output terminals of all programmable devices to one location, where they can be conveniently interconnected by various patch cords and plugs to meet the requirements of a given problem. Prewired problem boards can be exchanged at the patch bay in a few seconds and new coefficients set up typically in less than a half hour. Extensive automatic electronic patching systems have been developed to permit fast setup, as well as remote and time-shared operation. The analog computer basically represents an instrumentation of calculus, in that it is designed to solve ordinary differential equations. This capability lends itself to the implementation of simulated models of dynamic systems. The computer operates by generating voltages that behave like the physical or mathematical variables in the system under study. Each variable is represented as a continuously varying (or steady) voltage signal at the output of a programmed computational unit. Specific to the analog
computer is the fact that individual circuits are used for each feature or equation being represented, so that all variables are generated simultaneously. Thus the analog computer is a parallel computer in which the configuration of the computational units allows direct interactions of the computed variables at all times during the solution of a problem. Unique features. The unique features of the analog computer which are of value in science and technology are as follows: 1. Within the useful frequency bandwidth of the computational units and components, all programmed computations take place in parallel and are for practical purposes instantaneous. That is, there is no finite execution time associated with each mathematical operation, as is encountered with digital computer methods. 2. The dynamics of an analog model can be programmed for time expansion (slower than real system time), synchronous time (real time), or time compression (faster than real time). 3. The computer has a flexible addressing system so that almost every computed variable can be measured, viewed with vivid display instruments, and recorded at will. 4. One control mode of the computer, usually called HOLD, can freeze the dynamic response of a model to allow detailed examination of interrelationships at that instant in time, and then, after such study, the computing can be made to resume as though no stop had occurred. 5. By means of patch cords, plugs, switches, and adjustment knobs the analog computer program or model can be directly manipulated, especially during dynamic computation, and the resultant changes
Analog computer in responses observed and interpreted. 6. Because of the fast dynamic response of the analog computer, it is easy to implement implicit computation through the use of problem feedback. (This important and powerful mathematical quality of the analog computer is discussed more fully below.) 7. The computer can be used for on-line model building; that is, a computer model can be constructed in a step-by-step fashion directly at the console by interconnecting computational units on the basis of one-for-one analog representation of the real system elements. Then, by adjusting signal gains and attenuation parameters, dynamic behavior can be generated that corresponds to the desired response or is recognizable as that of the real system. This method allows a skillful person to create models when no rigorous mathematical equations for a system exist. 8. For those applications to which it is well suited, the analog computer operates at relatively low cost, thus affording the analyst ample opportunity to investigate, develop, and experiment within a broad range of parameters and functions. Components. Manipulations of the signals (voltages) in the analog computer are based upon the fundamental electrical properties associated with circuit components. The interrelation of voltages for representing mathematical functions is derived by combining currents at circuit nodes or junctions. See KIRCHHOFF’S LAWS OF ELECTRIC CIRCUITS; OHM’S LAW. The simplest arrangement of components for executing addition would be to impress the voltages to be added across individual resistors (Fig. 4a). The resistors would then be joined (Fig. 4b) to allow the currents to combine and to develop the output volt-
E1
R1
E3
E2 I2
l1
R3
R2
l3
(a)
E1
E2
E3
R1
summing junction IL
l1 R2
I2 I0
l3
R3
R0
E0
+
R0 I2 R2
RL
load
E0 = R0 I0 (b)
=
R0 l1 R1
+
R0 l R3 3
Fig. 4. Addition of electric currents by using a passive network of resistors. (a) Individual voltages and resistors. (b) Voltages and corresponding currents summed into a common resistor.
−A
Es
Eo = − A Es
Fig. 5. Symbol for a high-gain direct-current amplifier, with one inverting input referenced to ground.
age across a final resistor. Use of this simple configuration of elements for computation is impractical, because of the undesirable interaction between inputs. A change in one input signal (voltage) causes a change in the current that flows through the input resistor; this changes the voltage at the input resistor junction, and the change secondarily causes a different current to flow in the other input resistors. The situation gets more interactive when another computing circuit is attached so that part of the summing current flows away from the summing junction. This interaction effect also prevents exact computing. If, in some way, each voltage to be summed could be made independent of the other voltages connected to the summing junction, and if the required current fed to other circuits could be obtained without loading the summing junction, then precise computation would be possible. See DIRECT-CURRENT CIRCUIT THEORY. The electronic analog computer satisfies these needs by using high-gain (high-amplification) dc operational amplifiers. A symbol to represent a dc direct-coupled amplifier is shown in Fig. 5. According to convention, the rounded side represents the input to the amplifier, and the pointed end represents the amplifier output. A common reference level or ground exists between the amplifier input and output, and all voltages are measured with respect to it. The ground reference line is understood and is usually omitted from the symbol. The signal input network (consisting of summing resistors) connects to the inverting (or negative) amplifier input terminal. The noninverting (or positive) amplifier input terminal is normally connected to the reference ground. Generally the inverting input is called the summing junction (SJ) of the amplifier. Internal design of the amplifier is such that, if the signal at the summing junction is positive with respect to ground, the amplifier output voltage is negative with respect to ground. The amplifier has an open-loop voltage gain of −A; therefore an input voltage of Es results in an output voltage of −AEs. Gain of a commercial computing amplifier is typically 108; thus, an input voltage of less than 1 µV can produce several volts at the amplifier output. See AMPLIFIER; DIRECT-COUPLED AMPLIFIER. Because the operational amplifier thus inverts the signal (changes its sign or polarity), it lends itself to the use of negative feedback, whereby a portion of the output signal is returned to the input. This arrangement has the effect of lowering the net gain of the amplifier signal channel and of improving
625
626
Analog computer summing junction (SJ)
Rf
R in
E in
−A
Es
Eo
Fig. 6. High-gain amplifier which has been made into an operational amplifier through the inclusion of an input resistor and a feedback resistor tied together at the amplifier’s summing junction (SJ).
Rf ls = 0 E in
l in R in
lf SJ
Es 0
−A
Eo
l in = l f E in R in
= −
Eo = −
(1)
Eo
(2)
Rf
Rf E R in in
(3)
Fig. 7. Summing junction currents into an operational amplifier create the fixed gain function, determined by resistor values, as indicated in the box.
overall signal-to-noise ratio and increasing computational accuracy. Circuit operation (Fig. 6) can be viewed in the following manner. A dc voltage Ein applied to input resistor Rin produces a summing junction voltage Es. The voltage is amplified and appears at the amplifier output as voltage Eo (equal to −AEs, where A is voltage gain of the amplifier). Part of output voltage Eo returns through feedback resistor Rf to the summing junction. Because the returned or feedback voltage is of opposite (negative) polarity to the initial voltage at the summing junction, it tends to reduce the magnitude of Es, resulting in an overall input-output relationship that may be expressed as Eq. (1). In fact, Eo −1 −A = −1 = Ein A+1 1 + 1/A
A > 108
(1)
the summing junction voltage Es is so small that it
E1 E2 E3
R1
Rf
l1 lf
R2
lL SJ
l2 R3
Eo
RL l3 Eo = −
Rf R R E 1 + f E2 + f E 3 R3 R2 R1
Fig. 8. Operational amplifier which has been made into a summer through the use of several input resistors connected to the summing junction. Output is equal to the inverted weighted sum of the inputs.
is considered to be at practically zero, a condition called virtual ground. To illustrate how the operational amplifier serves the needs of the computing network, consider the currents that flow and the corresponding algebraic expressions (Fig. 7). The operational amplifier is designed to have high input impedance (high resistance to the flow of current into or out of its input terminal); consequently the amplifier input current Is can then be considered to be practically zero. The resulting current equation states that input current Iin is equal to feedback current If. Since the amplifier has a very high gain, the summing junction voltage is virtually zero. Voltage drop across Rin is thus equal to Ein; voltage drop across Rf is Eo. The equation in the box is the fundamental relationship for the amplifier. As long as the amplifier has such a high gain and requires a negligible current from the summing junction, the amplifier input and output voltages are related by the ratio of the two resistors and are thus not affected by the actual electronic construction of the amplifier. If several input resistors are connected to the same summing junction and voltages are applied to them (Fig. 8), then because the summing junction remains at practically zero potential, none of the inputs will interfere with the other inputs. Thus all inputs will exert independent and additive effects on the output. Because amplifier gain has a negative sign, output voltage equals the negative sum of the input voltages, weighted according to the values of the individual resistors in relation to the feedback resistor, as shown in the box in Fig. 8. When a computing circuit is connected to the amplifier output, a demand for load current is introduced. Such a required output load current must be supplied by the amplifier without a measurable change in its output voltage. That is, the operational amplifier must act as a voltage controller, supplying whatever current is required within limits, while maintaining the net voltage gain established by the mathematical ratio of its input and feedback elements. The operational amplifier and network is a linear network because once the input-feedback ratios are adjusted, signal channel gains remain constant (a straight-line function) during computation. Programming. To solve a problem using an analog computer, the problem solver goes through a procedure of general analysis, data preparation, analog circuit development, and patchboard programming. He or she may also make test runs of subprograms to examine partial-system dynamic responses before eventually running the full program to derive specific and final answers. The problem-solving procedure typically involves eight major steps, listed below: 1. The problem under study is described with a set of mathematical equations or, when that is not possible, the system configuration and the interrelations of component influences are defined in blockdiagram form, with each block described in terms of black-box input-output relationships. 2. Where necessary, the description of the system (equations or system block diagram) is rearranged
Analog computer
X
analog
digital
parallel machine
sequential machine
X Y
Y
analog-to-digital converter data transmission linkage
Z
digital-to-analog converter start
start 1 2 3 4 5
add multiply subtract store jump
data programs storage disk
finish
finish logic linkage monitoring and control linkage
scope
multichannel recorder
interactive terminal
run no. line printer
Fig. 9. Block diagram representation of a hybrid analog-to-digital computer.
in a form that may better suit the capabilities of the computer, that is, avoiding duplications or excessive numbers of computational units, or avoiding algebraic (nonintegrational) loops. 3. The assembled information is used to sketch out an analog circuit diagram which shows in detail how the computer could be programmed to handle the problem and achieve the objectives of the study. 4. System variables and parameters are then scaled to fall within the operational ranges of the computer. This may require revisions of the analog circuit diagram and choice of computational units. 5. The finalized circuit arrangement is patched on the computer problem board. 6. Numerical values are set up on the attenuators, the initial conditions of the entire system model established, and test values checked. 7. The computer is run to solve the equations or simulate the black boxes so that the resultant values or system responses can be obtained. This gives the initial answers and the “feel” for the system. 8. Multiple runs are made to check the responses for specific sets of parameters and to explore the influences of problem (system) changes, as well as the behavior which results when the system configuration is driven with different forcing functions. Hybrid computers. The accuracy of the calculations on a digital computer can often be increased through double precision techniques and more precise algorithms, but at the expense of extended solution time, due to the computer’s serial nature of operation. Also, the more computational steps there are to be done, the longer the digital computer will take to do them. On the other hand, the basic solution speed is very rapid on the analog computer because of its parallel nature, but increasing problem complexity demands larger computer size. Thus, for the analog computer the time remains the same regardless of
627
the complexity of the problem, but the size of the computer required grows with the problem. Interaction between the user and the computer during the course of any calculation, with the ability to vary parameters during computer runs, is a highly desirable and insight-generating part of computer usage. This hands-on interaction with the computed responses is simple to achieve with analog computers. For digital computers, interaction usually takes place through a computer keyboard terminal, between runs, or in an on-line stop-go mode. An often-utilized system combines the speed and interaction possibilities of an analog computer with the accuracy and programming flexibility of a digital computer. This combination is specifically designed into the hybrid computer. A modern analog-hybrid console (Fig. 2) contains the typical analog components plus a second patchboard area to include programmable, parallel logic circuits using high-speed gates for the functions of AND, NAND, OR, and NOR, as well as flip-flops, registers, and counters. The mode switches in the integrators are interfaced with the digital computer to permit fast iterations of dynamic runs under digital computer control. Data flow in many ways and formats between the analog computer with its fast, parallel circuits and the digital computer with its sequential, logic-controlled program (Fig. 9). Special high-speed analog-to-digital and digital-to-analog converters translate between the continuous signal representations of variables in the analog domain and the numerical representations of the digital computer. Control and logic signals are more directly compatible and require only level and timing compatibility. See ANALOG-TO-DIGITAL CONVERTER; BOOLEAN ALGEBRA; DIGITAL-TO-ANALOG CONVERTER. Programming of hybrid models is more complex than described above, requiring the user to consider
628
Analog states the parallel action of the analog computer interlaced with the step-by-step computations progression in the digital computer. For example, in simulating a space mission, the capsule control dynamics will typically be handled on the analog computer in continuous form, but interfaced with the digital computer, where the navigational trajectory is calculated. See COMPUTER. Per A. Holst Bibliography. R. M. Howe, Design Fundamentals of Analog Computer Components, 1961; G. A. Korn, Electronic Analog and Hybrid Computers, 1971; D. M. MacKay et al., Analog Computing at Ultra-High Speed, 1962; R. Tomovic, Repetitive Analog Computers, 1958.
Analog states States in neighboring nuclear isobars that have the same total angular momentum, parity, and isotopic spin. They also have nearly identical nuclear structure wave functions except for the transformation of one or more neutrons into an equivalent number of protons, which occupy the same single-particle states as the neutrons. Analog states (or isobaric analog states, IAS) have been observed throughout the periodic table, indicating that isotopic spin is a good quantum number. See ANGULAR MOMENTUM; I-SPIN; ISOBAR (NUCLEAR PHYSICS); NUCLEAR STRUCTURE; PARITY (QUANTUM MECHANICS). Since the nucleon–nucleon interaction has been found to be approximately charge-independent, it is possible to consider protons and neutrons as representing different charge states of a single particle, that is, a nucleon. Thus, a level (commonly referred to as a parent state) in a nucleus with Z protons and N neutrons can be expected to have an analog in the neighboring isobar with Z + 1 protons and N − 1 neutrons (and the same total number of nucleons, A = Z + N), where the protons and neutrons occupy the same orbits as those in the parent state. The energy difference between the parent and analog states predominantly arises from the increased contribution from the electrostatic Coulomb interaction to the total energy arising from the extra proton in the analog state. From this amount must be subtracted the neutron-proton mass difference of 0.782 MeV (energies are given on the atomic mass scale). The agreement between such calculated energies of analog states and their measured values is in general fairly precise but not exact. The reason is that small additional factors influence the level energies, such as electromagnetic effects, a small chargedependent nuclear interaction, isospin mixing, and nuclear structure effects. Isotopic spin. Since the nucleon has spin s = 1/2 , it obeys Fermi statistics. Charge independence of nuclear forces then implies that systems of two nucleons that can be described by the same wave functions that are antisymmetric in both space and spin would have identical energies. There are three such combinations of two-nucleon systems that can satisfy these conditions, namely, neutron-neutron, proton-
proton, and neutron-proton. The isotopic spin (also called isobaric spin or isospin) quantum number t = 1/2 has been introduced to distinguish between the two states of the nucleon, and provides a means to further identify these three states. The neutron is associated with a projection of t along the direction of quantization, that is, the z axis, with tz = 1/2 , and the proton with tz = −1/2 . The total z-projection of isospin, Tz, is given by the sum of the tz of the individual nucleons. Hence, Tz = 1, 0, −1 for the neutronneutron, neutron-proton, and proton-proton systems, respectively, and these represent an isobaric triplet with each system having total isospin T = 1. Only one two-nucleon system exists with a wave function that is symmetric in space and spin, that is, the neutron-proton combination for which T = TZ = 0. For a nucleus with N neutrons and Z protons, Tz = 1/2 (N − Z ) and, as with angular momentum, the value of isospin T must be equal to or greater than its z-projection, Tz. Analog states that differ solely in the interchange of either one neutron or proton (neighboring isobars) have values of Tz that differ by ±1. Those that differ by the interchange of two neutrons or protons have Tz that differ by ±2, and so forth. See EXCLUSION PRINCIPLE; FERMI-DIRAC STATISTICS. Nuclei with N = Z. Nuclei with N = Z (equal numbers of neutrons and protons) have Tz = 0. Hence, there can be no analog for states in such nuclei that also have total isospin T = 0. This prohibition follows because, as noted above, Tz changes by ±1 between neighboring isobars which must then have a total isospin T equal to or greater than 1. As a result, there is no analog of the ground state for these nuclei. However, excited states in N = Z nuclei for which T is greater than zero can have analogs. Nuclei with N > Z. Nuclei with N > Z (neutron excess) have Tz = 1/2 (N − Z), which is equal to or greater than 1/2 . Therefore, all states in these nuclei have T ≥ 1/2 , and in principle, they can thus have analogs. For heavy nuclei where the Coulomb energy becomes quite large, a question arose as to whether this large energy might destroy the isobaric symmetry. However, this question was resolved in 1961 with the observation of analog states in heavy nuclei by experiments that utilized the ( p, n) reaction. See NUCLEAR REACTION. Coulomb displacement energy. The energy required to transform a neutron into a proton in a nucleus with Z protons and N neutrons is called the Coulomb displacement energy, Ec. This energy can be calculated in terms of the energies of the parent and analog states from Eq. (1), where MZ+1 is the energy of the Ec = MZ+1 − MZ + np
(1)
analog state, MZ the energy of the parent state, and np = 0.782 MeV (as noted above) is the neutronproton mass difference (all energies are given on the atomic mass scale). The difference in the analog energies, that is, MZ+1 − MZ, can be expressed in terms of the Q-value for beta decay between the ground states of the isobars, and the energy of excitation of the analog state in the Z + 1 nucleus. From analyses of the abundant data on analog states
Analog states
18.826
∆Ec
T = 22 3.475 T = 22 3.198 T = 22 2.615
4− 5− 3−
T = 22 0.0
0+ 208 82Pb125
18.685 18.386 17.780
4− IAS T = 22 5− IAS T = 22 3− IAS T = 22
15.165
0+ IAS T = 22
4+ T = 21 5+ T = 21
0.063 0.0 208 83Bi125
Tz = 21 Qβ = 2.879
Tz = 22 Partial level diagram of the 208Pb and 208Bi isobar system. Broken lines connect parent–analog state pairs. Energies of excited states are given relative to the ground states, respectively, in megaelectronvolts. The isospin, T, and the angular momentum and parity, Jπ, are shown for each level.
throughout the periodic table, it is found that the Coulomb displacement energy can be represented ¯= to a high degree of accuracy by Eq. (2), where Z ¯ 1/3 ) − 0.861 (MeV) Ec = 1.412(Z/A
(2)
Z + 1/2 . The Coulomb displacement energy represents the minimum energy of excitation in the parent nucleus for which states with T = TZ + 1 can exist. Mass 208 system. A partial level diagram for the 208 Pb and 208Bi isobar system demonstrates the energy relationships discussed above (see illus.). In such a diagram, only a few of the many observed levels are shown. Below an excitation energy of Ec = 18.826 MeV in 208Pb, all levels have T = Tz = 22; the Coulomb displacement energy Ec is the minimum excitation at which T = 23 states can occur in 208Pb. The energy difference between the 208Pb and 208Bi ground states is Qβ = 2.879 MeV. The z-component of isospin for 208Bi is Tz = 21, which is also the total isospin, T, for all levels up to an excitation energy of 15.165 MeV. At this energy, a level (actually a resonance) is observed in 208Bi which possesses characteristics that are similar to those of the ground state of 208Pb, and is identified as the analog. The energy of excitation of the analog state lies in the nuclear continuum of the 208Bi mass system, that is, well above the energy at which 208Bi is unstable to particle decay. As a result, the analog state is accessible only by means of some nuclear reaction. It was first identified in low-energy studies of proton scattering on 207Pb, and later observed in the ( p, n) reaction on 208Pb. Schematically, in the ( p, n) reaction an energetic proton impinges upon a 208Pb nucleus and knocks out one of the excess neutrons that occupy the single-particle orbits that lie above those occupied by the protons, and the incident proton is captured into an equivalent orbit. This event
is known as a charge-exchange reaction (here, a neutron is transformed into a proton), and is related to the inverse reaction that occurs in positron beta decay. See RADIOACTIVITY. Fermi strength. Analog states are connected by either pure Fermi or mixed Fermi and Gamow-Teller transitions. The selection rules for Fermi transitions are J = 0, π = 0, T = 0, and Tz = ±1, where J represents the total angular momentum and π represents the parity. Those for Gamow-Teller transitions are J = 0 or ±1 (but transitions between analogs with J = 0 are forbidden), π = 0, T = 0 or ±1, and Tz = ±1. Both types of transitions require L = 0, where L represents the total orbital angular momentum. See SELECTION RULES (PHYSICS). Of special interest in the study of nuclear structure is the distribution of strength for various types of transitions based upon the ground state. For Fermi transitions, all of the strength is concentrated in the transition between the analog states. The total strength is proportional to the isospin matrix element between the analog states, which is given by Eq. (3), where |TTz represents the isospin part of the | T Tz − 1 |T− | T Tz |2 = (T + Tz )(T − Tz + 1) (3) parent wave function, and T− is an operator which changes a neutron into a proton. Since T = Tz, the total Fermi strength is proportional to (N − Z ), which is just equal to the number of excess neutrons. This outcome reflects the fact that each of the excess neutrons occupies an orbit (with specific space and spin descriptions) that lies at an energy above those occupied by the protons, and thus there is equal probability for each excess neutron to be transformed into a proton and to participate in the Fermi transition. Since the strength is well defined in practice, Fermi transitions can be used to obtain information pertaining to the reaction mechanisms that effect them. For beta decay, this approach usually involves determination of the vector coupling constant, GV, as well as fine details about the nuclear wave functions. In charge-exchange studies, knowledge is sought pertaining to details of the effective nucleon–nucleus interaction and, to some extent, details of the overlap between the spatial wave functions of the parent and analog states. Other applications. The study of analog states provides important information used to test nuclear theories. For example, the double charge-exchange reactions (π +, π −) [where π + and π − represent a pion with positive and negative charge, respectively] have been used to identify double isobaric analog states, that is, analogs in isobars removed by 2 charge units with Tz = T − 2. Such data have been useful for testing various formulas for predicting the relative masses of isobaric multiplets. Single and double charge-exchange reactions utilizing incident pions have also been used to investigate giant resonances built upon analog states. The single-particle structure of parent states can be studied by
629
630
Analog-to-digital converter observing the particle decay of the analog state when the decay resides in the nuclear continuum. Measurements of the widths of analog states provide information pertaining to their fragmentation, for example, their mixing with states of the same spin and parity but with total isospin lower by one unit, that is, equal to T − 1. See GIANT NUCLEAR RESONANCES. Daniel J. Horen Bibliography. M. S. Antony, J. Britz, and A. Pape, Coulomb displacement energies between analog levels for 44 ≤ A ≤ 239, Atom. Nucl. Data Tables, 40:9–56, 1988; M. S. Antony et al., Isobaric mass equation for A = 1–45 and systematics of Coulomb displacement energies, Atom. Nucl. Data Tables, 33:447–478, 1985; A. Bohr and B. R. Mottelson, Nuclear Structure, vol. 1, 1969, vol. 2, 1975; D. H. Wilkinson (ed.), Isospin in Nuclear Physics, 1969.
Analog-to-digital converter A device for converting the information contained in the value or magnitude of some characteristics of an input signal, compared to a standard or reference, to information in the form of discrete states of a signal, usually with numerical values assigned to the various combinations of discrete states of the signal. Analog-to-digital (A/D) converters are used to translate analog information, such as audio signals or measurements of physical variables (for example, temperature, force, or shaft rotation), into a form suitable for digital handling, which might involve any of these operations: (1) processing by a computer or by logic circuits, including arithmetical operations, comparison, sorting, ordering, and code conversion, (2) storage until ready for further handling, (3) dis000
001 010 011 100 101 110
0 8
1 8
2 8
3 8
4 8
5 8
6 8
111
7 8
8 8
binary code analog input range (full scale)
Fig. 1. A three-bit binary representation of a range of input signals.
start conversion
"busy"
analog signal A/D converter ("ADC")
bit 1 (MSB) bit 2 bit 3 bit 4
reference
bit 5 (LSB)
clock serial output 5V ("1")
time
16 0 4 2 0 + + + + 32 32 32 32 32
0V ("0") 5V ("1")
play in numerical or graphical form, and (4) transmission. If a wide-range analog signal can be sampled and converted, with sufficient frequency, to an appropriate number of two-level digits, or bits, to represent each sample or snapshot, the digital representation of the signal can be transmitted through a noisy medium without essential degradation of the fine structure of the original signal. See COMPUTER GRAPHICS; DATA COMMUNICATIONS; DIGITAL COMPUTER. Concepts and structure. Conversion involves quantizing and encoding. Quantizing means partitioning the analog signal range into a number of discrete quanta and determining to which quantum the input signal belongs. Encoding means assigning a unique digital code to each quantum and determining the code that corresponds to the input signal. The most common system is binary, in which there are 2n quanta (where n is some whole number), numbered consecutively; the code is a set of n physical twovalued levels or bits (1 or 0) corresponding to the binary number associated with the signal quantum. [The term binary is used here in two senses: a binary (radix-2) numbering system, and the binary (twovalued) decisions forming each bit.] See BIT. Figure 1 shows a typical three-bit binary representation of a range of input signals, partitioned into eight (that is, 23) quanta. For example, a signal in the vicinity of 3/8 full scale (between 5/16 and 7/16) will be coded 011 (binary 3). See NUMBERING SYSTEMS. Conceptually, the conversion can be made to take place in any kind of medium: electrical, mechanical, fluid, optical, and so on (for example, shaft-rotationto-optical); but by far the most commonly employed form of A/D converters comprises those devices that convert electrical voltages or currents to coded sets of binary electrical levels (for example, +5 V or 0 V) in simultaneous (parallel) or pulse-train (serial) form, as shown in Fig. 2. Serial outputs are not always made available in the form of binary numbers. The converter depicted in Fig. 2 converts the analog input to a five-digit “word.” If the coding is binary, the first digit (most significant bit, abbreviated MSB) has a weight of 1/2 full scale, the second 1/4 full scale, and so on, down to the nth digit (leastsignificant bit, abbreviated LSB), which has a weight of 2−n of full scale (1/32 in this example). Thus, for the output word shown, 10110, the analog input must be given approximately by the equation below.
parallel outputs
5V ("1") 0V ("0")
(analog) (digital) power-supply leads Fig. 2. Basic diagram of an analog-to-digital converter, showing parallel and serial (return-to-zero) output formats for code 10110.
=
11 22 = FS (full scale) 32 16
The number of bits, n, characterizes the resolution of a converter. The table translates bits into other conventional measures of resolution in a binary system. Figure 2 also shows a commonly used configuration of connections to an A/D converter: the analog signal and reference inputs; the parallel and serial
Analog-to-digital converter
631
Binary resolution equivalents∗ 2⫺n
Bit †
FS MSB‡ 2 3 4 5 6 7 8 9 10 11 12 13 14 15 16 17 18 19 20 ∗
0
2 2⫺1 2⫺2 2⫺3 2⫺4 2⫺5 2⫺6 2⫺7 2⫺8 2⫺9 2⫺10 2⫺11 2⫺12 2⫺13 2⫺14 2⫺15 2⫺16 2⫺17 2⫺18 2⫺19 2⫺20
1/ 2n (fraction)
dB (decibels)
1/ 2n (decimal)
1 1/ 2 1/ 4 1/ 8 1/ 16 1/ 32 1/ 64 1/ 128 1/ 256 1/ 512 1/ 1024 1/ 2048 1/ 4096 1/ 8192 1/ 16, 384 1/ 32, 768 1/ 65, 536 1/ 131, 072 1/ 262, 144 1/ 524, 288 1/ 1, 048, 576
0 ⫺6 ⫺12 ⫺18.1 ⫺24.1 ⫺30.1 ⫺36.1 ⫺42.1 ⫺48.2 ⫺54.2 ⫺60.2 ⫺66.2 ⫺72.2 ⫺78.3 ⫺84.3 ⫺90.3 ⫺96.3 ⫺102.3 ⫺108.4 ⫺114.4 ⫺120.4
1.0 0.5 0.25 0.125 0.0625 0.03125 0.015625 0.007812 0.003906 0.001953 0.0009766 0.00048828 0.00024414 0.00012207 0.000061035 0.0000305176 0.0000152588 0.00000762939 0.000003814697 0.000001907349 0.0000009536743
%
Parts per million
100 50 25 12.5 6.2 3.1 1.6 0.8 0.4 0.2 0.1 0.05 0.024 0.012 0.006 0.003 0.0015 0.0008 0.0004 0.0002 0.0001
1,000,000 500,000 250,000 125,000 62,500 31,250 15,625 7,812 3,906 1,953 977 488 244 122 61 31 15 7.6 3.8 1.9 0.95
From D. H. Sheingold (ed.), Analog-Digital Conversion Handbook, 3d ed., Analog Devices, Inc., 1986.
† Full scale. ‡ Most significant
bit.
digital outputs; the leads from the power supply, which provides the required energy for operation; and two control leads—a start-conversion input and a status-indicating output (busy) that indicates when a conversion is in progress. Serial-output converters often furnish a clock output pulse to indicate the time at which a bit is available. The reference voltage or current is often developed within the converter. Second in importance to the binary code and its many variations is the binary-coded decimal (BCD), which is used when the encoded material is to be displayed or printed in numerical form. In BCD, each digit of a radix-10 number is represented by a four-digit binary subgroup. For example, the BCD code for 379 is 0011 0111 1001. The output of the A/D converter used in digital panel meters is usually BCD. Techniques. There are many techniques used for A/D conversion, ranging from simple voltage-level comparators to sophisticated closed-loop systems, depending on the input level, output format, control features, and the desired speed, resolution, and accuracy. The three most popular techniques, used individually or in combination, are direct or flash conversion; integrating conversion, exemplified by dual slope; and feedback conversion, exemplified by the method of successive approximations. Comparator. The most elementary A/D converter is the comparator, a 1-bit (2-value) converter. It has two inputs, the signal and the reference: whenever the signal is greater than the reference, the output is 1; when less, the output is 0. A digitally controlled latch may be included to lock in the value of the comparison at a definite time. Direct conversion. The fastest, and conceptually simplest, approach to n-bit conversion is via the direct or flash converter. It quantizes by using 2n − 1 com-
parators to compare the input signal simultaneously to a set of reference signals representing all amplitude quanta. Then it employs a logical scheme called a priority encoder to determine the parallel digital output corresponding to a given set of comparator output states. Figure 3 is a simplified block diagram of a 2-bit (4-level) flash converter. The resistor string acts as a
reference comparators 4.0 V
2.5 V
−
nominal input 3V 2V 1V 0V
quantum
code
111 011
11 10
001 000
01 00
+ "1" for V in > 2.5 V resistive divider
1.5 V
− MSB
+ "1" for V in > 1.5 V 0.5 V
−
LSB
+ "1" for V in > 0.5 V
input (V in ) Fig. 3. Diagram of a 2-bit flash converter.
priority encoder
632
Analog-to-digital converter
analog input, V in comparator
integrator opposite polarity reference, V ref
output voltage
start conversion
control logic
1 T
overrange
counter MSB
(a)
−
switch from V in to V ref
∫ V in dt
1 T
∫ V ref dt stop when
K T
V in N0 −
K T
V ref N = 0
time
start integrate V in for N 0 counts
LSB digital output at end of conversion
∫ V in dt = KT V in N 0
1 T
clock
"busy"
integrate V ref (N counts)
v in N = N0 v ref
(b)
Fig. 4. Example of a dual-slope conversion. (a) Block diagram of converter. (b) Integrator output. Here, k is a constant, and T is the RC (resistor-capacitor) time constant of the integrator.
voltage divider to provide reference voltages for the three (that is, 22 − 1) comparators. When the common input is below the lowest reference level, all of the comparators have 0 output. When the input is greater than the reference for the lowest comparator, its output switches to 1; when the next comparator’s reference is exceeded, its output too switches to 1; and when the third comparator’s reference level is exceeded, all three comparators are at 1. Thus, the length of the series of 1’s identifies the input voltage’s quantum. The quantum must now be encoded as a 2-bit digital word at the output. In the priority-encoder scheme shown in Fig. 3, the output is a binary word, and the correspondence between quantum and code is given in the table.
Practical flash converters require lengthy resistor strings, large numbers of comparators, and rather complicated encoders. For example, a 10-bit flash converter calls for 1023 comparators and requires encoding 1024 levels into a 10-bit word. Nevertheless, in integrated-circuit form, all of the circuity exists on a single monolithic silicon chip. Integrating conversion. Unlike flash converters, which measure the input instant by instant, integrating converters use analog integrating circuits in a variety of circuit architectures to convert average values of the input into trains of pulses, which are then encoded by digital counters or processors. Integrating converters have high resolution and low noise sensitivity, and the train of pulses can be transmitted as a high-level noise-insensitive signal for encoding output
analog input
comparator
analog value of digital output reference
analog input + 1/2 LSB
digital-to-analog converter
parallel digital output
(a)
accept
16 32
LSB
start conversion "busy"
reject
20 32
MSB
serial output
reject
24 32
sequencing register
clock
accept
accept MSB
12 32
LSB time
1
0
1
1
0
control logic (b) Fig. 5. Successive-approximations conversion. (a) Block diagram of converter. (b) Digital-to-analog converter output for the example in Fig. 2.
Analysis of variance at a remote location. The class includes voltage-tofrequency converters, delta-modulation types, and dual-slope converters. The dual-slope converter, discussed here as an example of the class, has long been popular in digital voltmeters; it can perform complete conversions with high accuracy at low speeds, generally a few conversions per second. Figure 4a is a simplified block diagram of a dual-slope converter. The input is integrated for a period of time determined by a clock-pulse generator and counter (Fig. 4b). The final value of the signal integral becomes the initial condition for integration of the reference in the opposite sense, while the clock output is counted. When the net integral is zero, the count stops. Since the integral “up” of the input over a fixed time (N0 counts) is equal to the integral “down” of the fixed reference, the ratio of the number of counts of the variable period to that of the fixed period is equal to the ratio of the average value of the signal to the reference. Feedback conversion. Some A/D converter types employ an n-bit digital-to-analog converter to compare an analog version of the digital output with the analog input; the decision (higher or lower) causes a change in the digital output to bring its value closer to the input value. The change may be made by a counter that increases or decreases the digital value in small steps, or by the educated guessing of a successive-approximations converter. Successive-approximations conversion combines superior speed and accuracy; it is the principal technique used for data-acquisition and computerinterface systems. Figure 5a is a simplified block diagram of a successive-approximations converter. In a manner analogous to the operation of an apothecary’s scale with a set of binary weights, the input is weighed against a set of successively smaller fractions of the reference, produced by a digitalto-analog (D/A) converter that reflects the number in the output register. See DIGITAL-TO-ANALOG CONVERTER. First, the MSB is tried (1/2 full scale). If the signal is less than the MSB, the MSB code is returned to zero; if the signal is equal to or greater than the MSB, the MSB code is latched in the output register (Fig. 5b). The second bit is tried (1/4 full scale). If the signal is less than 1/4 or 3/4, depending on the previous choice, bit 2 is set to zero; if the signal is equal to or greater than 1/4 or 3/4, bit 2 is retained in the output register. The third bit is tried (1/8 full scale). If the signal is less than 1/8, 3/8, 5/8, or 7/8, depending on previous choices, bit 2 is set to zero; otherwise, it is accepted. The trial continues until the contribution of the least-significant bit has been weighed and either accepted or rejected. The conversion is then complete. The digital code latched in the output register is the digital equivalent of the analog input signal. Physical electronics. The earliest A/D converters were large rack-panel chassis-type modules using vacuum tubes, requiring about 1.4 ft3 (0.04 m3) of space and many watts of power. Since then,
they have become smaller in size and cost, evolving through circuit-board, encapsulated-module, and hybrid construction, with improved speed and resolution. Single-chip 12-bit A/D converters with the ability to interface with microprocessors are now available in small integrated-circuit packages. Integrated-circuit A/D converters, with 8-bit and better resolution and conversion rates of hundreds of megahertz, are also commercially available. See MICROPROCESSOR. Daniel H. Sheingold Bibliography. Analog Devices, Inc. Engineering Staff, Analog-Digital Conversion Handbook, ed. by D. H. Sheingold, 3d ed., 1986; M. Demler, HighSpeed Analog to Digital Conversion, 1991; D. F. Hoeschele, Jr., Analog-to-Digital and Digital-toAnalog Conversion Techniques, 2d ed., 1994; R. Van de Plassche, Integrated Analog-to-Digital and Digital-to-Analog Converters, 1994.
Analysis of variance Total variation in experimental data is partitioned into components assignable to specific sources by the analysis of variance. This statistical technique is applicable to data for which (1) effects of sources are additive, (2) uncontrolled or unexplained experimental variations (which are grouped as experimental errors) are independent of other sources of variation, (3) variance of experimental errors is homogeneous, and (4) experimental errors follow a normal distribution. When data depart from these assumptions, one must exercise extreme care in interpreting the results of an analysis of variance. Statistical tests indicate the contribution of the components to the observed variation. In an illustrative experiment, t methods of treatment are under study, and n samples are measured for each treatment for a total of nt samples. Measurement Xij of the ith sample that received the jth treatment records an overall effect µ, an effect β j produced by the jth treatment, and an effect ij produced by experimental error. The three effects are additive, so that Eq. (1) holds, where i = 1, . . . , n; Xij = µ + β j + ij
(1)
and j = 1, . . . , t. The statistical problem is to test for the existence of these effects. The analysis of variance in this example is presented in the table. Entries in the sum of squares column represent that part of the total variation that is attributable to each source. Total sum of squares Q is the sum over all squared deviations of observations Xj from the grand mean X, Eq. (2). Similarly, X = ( i Xij)/nt
(2)
within treatments, sum of squares E is the sum over all squared deviations of observations Xij within a treatment from the mean X j of that treatment, Eq. (3). X j = ( i Xij)/n
(3)
Also, between treatments, sum of squares T is n times
633
634
Analytic geometry Illustrative example of the analysis of variance Source of variation
Sum of squares
Degrees of freedom
Between treatments
T ⫽ nΣ j (X j ⫺ X)2
t⫺1
Within treatments
E ⫽ Σ i,j (X ij ⫺ X j )2
t(n ⫺ 1)
Total
Q ⫽ Σ i,j (X ij ⫺ X)2
nt ⫺ 1
the sum over all treatments of the squared deviations of treatment means X j from a grand mean X as defined by Eqs. (2) and (3). The sum of squares is generally computed more easily from the equivalent formulas (4)–(6). Q = i,jXij2 − T = j
( i,jXij)2 nt
( i Xij)2 ( i, jXij)2 − n nt E =Q−T
(4)
Mean square T ⫽ T / (t ⫺ 1) E ′ ⫽ E / t(n ⫺ 1)
at the α point on the F distribution with t − 1 and t(n − 1) degrees of freedom. See BIOMETRICS; QUALITY CONTROL; STATISTICS. Robert L. Brickley Bibliography. D. R. Anderson, D. F. Sweeny, and T. A. Williams, Introduction to Statistics: Concepts and Applications, 3d ed., 1993; S. Jarrell, Basic Statistics, 1994; N. A. Weiss, Introductory Statistics, 5th ed., 1999.
(5) (6)
The entries under degrees of freedom represent the number of independent comparisions upon which the sum of squares for the source of variation is based. In every case the linear restriction imposed by the relationship of the particular mean to the observations results in the loss of one degree of freedom. Therefore the number of degrees of freedom is always one less than the number of deviations used to compute the sum of squares. The mean squares in the analysis of variance are obtained by dividing the sum of squares by the corresponding degrees of freedom. The within-treatments mean square is an estimate of σ 2, the variance of the error term ij in the additive model. It represents the random or unexplained variation in the data. The between-treatments mean square is an estimate of σ 2 + σ β2, where σ β2 is the variance of the treatment effects β j. If the treatment means differ substantially, the β j effects estimated by (X j − X) will differ correspondingly and will have a large variance σ β2. If on the other hand the means do not differ, the treatment effects β j would be zero and σ β2 would be zero. In this case the treatment mean square would be equal to the error mean square and both would be independent estimates of σ 2. By comparing the ratio T/E of between treatment mean square T to within-treatment mean square E with unity, the variation due to treatments is compared with the variation due to random or unexplained factors. If this ratio, called the F ratio, is close to unity, there is no evidence of a treatment effect. However, if ratio T/E is substantially greater than unity there may be a significant treatment effect. To compare the mean squares objectively, one uses the F test of significance in which the statistical hypothesis is that σ β2 = 0. Under this hypothesis it can be concluded that the treatment effects are significantly different from zero at the significance level α if the calculated F ratio is greater than the value of F
Analytic geometry A branch of mathematics in which algebra is applied to the study of geometry. Because algebraic methods were first systematically applied to geometry in 1637 by the French philosopher-mathematician Ren´e Descartes, the subject is also called cartesian geometry. The basis for an algebraic treatment of geometry is provided by the existence of a one-to-one correspondence between the elements, “points” of a directed line g, and the elements, “numbers,” that form the set of all real numbers. Such a correspondence establishes a coordinate system on g, and the number corresponding to a point of g is called its coordinate. The point O of g with coordinate zero is the origin of the coordinate system. A coordinate system on g is cartesian provided that for each point P of g, its coordinate is the directed distance OP. Then all points of g on one side of O have positive coordinates (forming the positive half of g) and all points on the other side have negative coordinates. The point with coordinate 1 is called the unit point. Since the relation OP + PQ = OQ is clearly valid for each two points P, Q, of directed line g, then PQ = OQ − OP = q − p, where p and q are the coordinates of P and Q, respectively. Those points of g between P and Q, together with P, Q, form a line segment. In analytic geometry it is convenient to direct segments, writing PQ or QP accordingly as the segment is directed from P to Q or from Q to P, respectively. To find the coordinate of the point P that divides the segment P1P2 in a given ratio r, put P 1P/P 2P = r. Then (x − x1)/(x − x2) = r, where x1, x2, x are the coordinates of P1, P2, P, respectively, and solving for x gives x = (x1 − rx2)/(1 − r). Clearly r is negative for each point between P1, P2 and is positive for each point of g external to the segment. The midpoint of the segment divides it in the ratio −1, and hence its coordinate x = (x1 + x2)/2. See MATHEMATICS. The plane. Choose any two intersecting lines g1, g2 of the plane, with cartesian coordinate systems selected on each so that the intersection point has coordinate O in each system. To each point P of the
Analytic geometry plane an ordered pair of numbers (x, y) is attached as coordinates, where x is the coordinate of the point of intersection of g1 with the line through P parallel to g2, and y is the coordinate of the point of intersection of g2 with the line through P parallel to g1. If P lies on g, its coordinates are (x, 0), where x is the coordinate of P in the coordinate system on g1. Similarly, each point of g2 has coordinates (0, y). The origin O has coordinates (0, 0). The lines g1 and g2 are called the x axis and y axis, respectively. It is usually convenient to take the same scale on each axis; that is, the segments joining the unit points on the two axes to the origin are congruent. The notation P(x, y) denotes a point P with coordinates (x, y). If ω denotes the angle made by the positive halves of the two axes, and d the distance of points P1(x1, y1), P2(x2, y2), application of the law of cosines yields Eq. (1). d = (x1 − x2 )2 + (y1 − y2 )2 1/2 + 2(x1 − x2 )(y1 − y2 ) cos ω (1)
of certain important figures among curves and surfaces. See CURVE FITTING. Equations of lines. A line may be determined by data of various kinds, each yielding a different form for its equation. Thus an equation for a line g through P1(x1, y1) perpendicular to the x axis is evidently x = x1. If g goes through P1(x1, y1), P2(x2, y2) and is not perpendicular to the x axis, let P(x, y) be the coordinates of any other point of g. Then (y − y1)/(x − x1) = λ = (y2 − y1)/(x2 − x1), whereas if (x, y) are the coordinates of a point not on g, (y − y1)/(x − x1) = λ. Hence an equation for g is shown as Eq. (4), from which follow Eqs. (5) and (6).
Since cos 90◦ = 0, this important formula is simplified by taking the axes mutually perpendicular. Though it is occasionally useful to employ oblique axes (that is, ω = 90◦), the simplifications resulting from a rectangular cartesian coordinate system make it the usual choice. Such a cartesian coordinate system is assumed in what follows. Thus, the distance d of P1(x1, y1), P2(x2, y2) is given by Eq. (2).
The two-point form is given by Eq. (4), the pointslope form by Eq. (5), and the symmetric form by Eq. (6). The validity of the determinant form, x y 1 x1 y1 1 = 0 x2 y2 1
d = [(x1 − x2 )2 + (y1 − y2 )2 ]1/2
(2)
Let θ denote the smaller of the two angles that a line g makes with the positive half of the x axis, measured in the (positive) direction of rotation (that is, from the positive half of the x axis to the positive half of the y axis). Angle θ is called the slope angle of g, and the number λ = tan θ, the slope of g, plays an important role in plane analytic geometry. Slope is not defined for any line perpendicular to the x axis. If P1(x1, y1), P2(x2, y2) are two distinct points of line g with slope angle θ = 90◦, and Q denotes the intersection of the line through P1 perpendicular to the y axis with the line through P2 perpendicular to the x axis, then Eq. (3) holds. λ = tan θ =
y2 − y1 QP2 = x P1 Q 2 − x1
(3)
Loci and equations. The correspondence between the geometric entity “point” and the arithmetic entity “pair of real numbers,” upon which plane analytic geometry is based, results in associating with each geometric locus one or more equations that are satisfied by the coordinates of all those (and only those) points forming the locus (equations of the locus), and in associating with each system of equations in the variables x, y the figure (graph of the equations) whose points are determined by the pairs of numbers satisfying the equations. Thus the algebraic method of studying geometry is balanced by a geometric interpretation of algebra. A central problem in analytic geometry is that of finding equations
y − y1 y2 − y1 = x − x1 x2 − x1
(4)
y − y1 = λ(x − x1 )
(5)
y − y1 x − x1 = x2 − x1 y2 − y 1
(6)
as well as the slope-intercept form, y = λx + b, and the intercept for x/a + y/b = 1, a · b = 0, also follow readily from Eq. (1). See DETERMINANT. A line g is determined by its distance p from the ori> 0) and the angle β that the perpendicular gin O ( p = to g from O makes with the x axis. (The perpendicular is directed from O to g in case g does not go through O, and so as to make β < 180◦ in the contrary case.) The equation of g in terms of p and β (the so-called normal form) is x cos β + y sin β − p = 0. The directed distance from g to any point P(x0, y0) is x0 cos β + y0 sin β − p. It follows that the equations of the bisectors of the angles formed by two lines x cos βi + y sin βi − pi = 0
(i = 1, 2)
are x cos β1 + y sin β1 − p1 = ±(x cos β2 + y sin β2 − p2 ) The general form of an equation of a line is Ax + By + C = 0, where A, B are not both zero. The normal form is obtained from the general √ form upon dividing the left-hand member by ± A2 + B2, choosing the sign of the radical opposite to that of C in case C = 0, the same as that of B in case C = 0 and B = 0, and the same as that of A in case B = 0 and C = 0. If B = 0, the general form reduces to x = constant, a line perpendicular to the x axis, and if B = 0, one obtains y = − (A/B)x + C/A, a line with slope − A/B and y intercept C/A. Thus to each line there corresponds a linear equation in x and y, and with each such equation is associated a line.
635
636
Analytic geometry Angle between two lines. The angle φ (0◦ < φ < 180◦) from a line g1 to another intersecting line g2 is that through which g1 must be rotated (about the point of intersection) to coincide with g2. If θ 1, θ 2 are the slope angles of g1, g2, respectively, then φ = θ2 − θ 1 and tan φ =
tan θ2 − tan θ1 λ 2 − λ1 = 1 + tan θ1 tan θ2 1 + λ 1 · λ2
provided θ 1 = 90◦ = θ 2. Consequently, g1 and g2 are mutually perpendicular if and only if 1 + λ1 · λ2 = 0. The formula for tan φ holds in case g1, g2 are parallel. Then φ = 0◦ and λ1 = λ2. It follows that two distinct lines Aix + Biy + Ci = 0, with i = 1, 2, are mutually perpendicular if and only if A1A2 + B1B2 = 0, and parallel provided A1B2 − A2B1 = 0. Area of a triangle. Let Pi(xi, yi), with i = 1, 2, 3, be vertices of a triangle whose area is denoted by A. Then A = 12 d · D, where d = [(x1 − x2)2 + (y1 − y2)2]1/2 and D is the distance of P3 from the line joining P1, P2. Substituting the coordinates (x3, y3) of P3 for (x, y) in the normal form of the equation of that line gives ± (y2 − y1 )x3 − (x2 − x1 )y3 − x1 y2 + x2 y1 D= (x2 − x1 )2 + (y2 − y1 )2 and hence A = ±(1/2 )(x1 y2 + x2 y3 + x3 y1 −x3 y2 − x2 y1 − x1 y3 ) or x1 A = ±(1/2) · x2 x3
y1 y2 y3
1 1 1
The positive sign holds provided the vertices P1, P2, P3 are in counterclockwise order, and the negative sign in the contrary case. Linear combinations. If u1 = 0, u2 = 0 are equations of lines through a point P, for every choice of constants c1, c2 (except c1 = c2 = 0) the linear combination c1u1 + c2u2 = 0 is an equation of a line through P, and every line through P has an equation of that form. It follows that three lines ui = 0, with i = 1, 2, 3, are concurrent provided there exist constants c1, c2, c3 (not all zero) such that the linear combination c1u1 + c2u2 + c3u3 = 0 for every pair of numbers (x, y). Putting ui = Ai x + Bi y + Ci, with i = 1, 2, 3, then c1, c2, c3 are nontrivial solutions of the system of equations c1A1 + c2A2 + c3A3 = 0, c1B1 + c2B2 + c3B3 = 0, c1C1 + c2C2 + c3C3 = 0. Hence the lines ui = 0, with i = 1, 2, 3, are concurrent if and only if A1 B1 C1 A2 B2 C2 = 0 A3 B3 C3
Circle. By use of the formula for the distance of two points and the definition of a circle, an equation for > 0) is a circle with center C(x0, y0) and radius r (r = found to be (x − x0)2 + (y − y0)2 = r2. If r = 0, the only (real) point of the locus is the center (x0, y0) and the circle is a point circle. The above equation is called the standard form. Expansion yields x2 + y2 − 2x0x − 2y0y − x02 + y02 − r2 = 0, a quadratic in x, y with equal, nonzero, coefficients of x2 and y2, and with the product term xy lacking. Conversely, the locus of each such equation A(x2 + y2) + 2Dx + 2Ey + F = 0, A = 0, is a (real) circle with center (−D/A, −E/A) and radius r = (1/A) · [D2 + E2 − AF]1/2 provided D2 + > 0, for by “completing the square” the E2 − AF = above equation may be put in the standard form. Let P(x1, y1) be on the circle (x − x0)2 + (y − y0)2 = r2. An equation for the tangent to the circle at P is easily seen to be (y1 − y0) (y − y1) + (x1 − x0) (x − x1) = 0. Adding (x1 − x0)2 + (y1 − y0)2 to the left side of this equation, and its equal r2 [since P(x1, y1) lies on the circle] to the right side, yields (x1 − x0)(x − x0) + (y1 − y0)(y − y0) = r2 as an equation of the tangent. The formal process by which this equation may be obtained from the standard form of the equation of a circle is called polarization. If P(x1, y1) is not on the circle, the line that is the locus of that equation is called the polar line of P with respect to the circle. If P is outside the circle, its polar line joins the contact points of the two tangents from P to the circle. See CIRCLE. Polarization of the general equation of a circle gives A(x1x + y1y) + D(x + x1) + E(y + y1) + F = 0, as an equation of the tangent at (x1, y1) if that point is on the circle, and of the polar line of (x1, y1) otherwise. The tangential distance t from (x1, y1) to the circle (x − x0)2 + (y − y0)2 = r2 is given by t2 = (x1 − x0)2 + (y1 − y0)2 − r2. Hence (x − x0)2 + (y − y0)2 − r02 = (x − x1)2 + (y − y1)2 − r12 is an equation of the locus of points with equal tangential distances from the two circles with centers (x0, y0), (x1, y1) and radii r0, r1, respectively. Since the equation is linear, it represents a line which evidently contains the common chord of the circles in case they intersect in two distinct points. If ui = 0, with i = 1, 2, 3, denotes equations of three circles that intersect pairwise in two distinct points, then the three common chords are concurrent, since the equation u1 − u2 = 0 is a linear combination of the equations u2 − u3 = 0, u3 − u1 = 0. An equation of the circle through three noncollinear points Pi(xi,yi), with i = 1, 2, 3, may be written 2 x + y2 x y 1 x2 + y2 x y 1 1 1 1 1 2 x2 + y22 x2 y2 1 x2 + y2 x y 1 3
3
3
3
Conic sections. Much of plane analytic geometry deals with a class of curves which (from the way in which they were first studied) are known as conic sections or conics. A conic is the locus of a point P that moves so that its distance from a fixed point F (the focus) is in a constant positive ratio (the
Analytic geometry eccentricity) to its distance from a fixed line (the directrix) not through F. Let (c, 0) be the coordinates of F, c > 0, and take the y axis as the directrix. Then P(x, y) satisfies the equation [(x − c)2 + y2]1/2 = x; that is, (1 − 2)x2 + y2 − 2cx + c2 = 0, and it is easily seen that each point whose coordinates satisfy this equation is on the conic. Hence each conic is represented by a second-degree equation in the cartesian coordinates (x, y). A conic is called a parabola, ellipse, or hyperbola accordingly as = 1, < 1, > 1, respectively. See CONIC SECTION. Parabola. When = 1, the equation for the conic obtained above becomes y2 − 2cx + c2 = 0 or y2 = 2c(x − c/2). The curve has an axis of symmetry (the x axis when the focus and directrix are chosen as above). The point (c/2, 0) at which the axis intersects the curve is the vertex. Putting x = c in the equation, it is seen that the chord through the focus that is perpendicular to the axis (the latus rectum) has length 2c. All quadratics, x = Ay2 + By + C, y = Ax2 + Bx + C, A = 0, are equations of parabolas (with axes perpendicular to the y axis, or the x axis, respectively). All other parabolas have equations Ax2 + 2Bxy + Cy2 + 2Dx + 2Ey + F = 0, with discriminant A B D B C E = 0 and AC − B2 = 0 D E F (A quadratic with nonvanishing discriminant is called irreducible.) A simple standard form of an equation of a parabola is obtained by taking (c/2, 0) for focus (c > 0) and the line x = − c/2 for directrix, resulting in y2 = 2cx. Polarizing gives y1y = c(x + x1) for an equation of the tangent at the point P(x1, y1) on the parabola. This tangent cuts the axis at the point Q(−x1, 0) whose distance from the focus is x1 + c/2. But this is the distance from the directrix to P, and consequently the triangle FPQ is isosceles with - FQP = - FPQ. It follows that the line joining F to any point P of the parabola makes the same angle with the tangent at P as does the line through P that is parallel to the axis of the parabola. Thus each light ray emanating from the focus is reflected parallel to the axis. See PARABOLA. Ellipse. When < 1, the equation of the general conic given above becomes x−c 2 y2 2 c2 + = ( < 1) 2 2 1− 1− (1 − 2 )2 )
)
This is an equation of the ellipse with focus (c, 0) and directrix the y axis. Putting X = x − c/(1 − 2), Y = y, the standard form X2/A2 + Y2/B2 = 1 of the equation is obtained, where A = c/(1 − 2) > 0, B = c/(1 − 2)1/2 > 0. This is equivalent to referring the ellipse to a new set of coordinate axes obtained by translating the origin to the point [c/(1 − 2),0]. The coordinates of the focus F in the XY coordinate system are [−c 2/(1 − 2),0] and the equation of the directrix is X = −c/(1 − 2). Since the standard form shows the curve to be symmetric with respect to each of the new axes and the origin, it is clear that F[c 2/(1 − 2),0] and X = c/(1 − 2) may be taken
as focus and directrix. Putting C = c 2/(1 − 2), then C2 = A2 − B2, and X 2/A2 + Y 2/B2 = 1 is seen to be the locus of points P(X,Y), the sum of whose distances from (C, 0) and (−C, 0) is 2A. This relation is frequently used to define an ellipse. The chord containing the foci is the major axis; its length is 2A. The chord of length 2B through the midpoint of the foci, perpendicular to the major axis, is the minor axis. The tangent to X2/A 2 + Y 2/B2 = 1 at P(X1,Y1) on the ellipse is X1X/A2 + Y1Y/B2 = 1. Lines PF, PF make equal angles with the tangent, so sound or light that emanates from one focus is reflected to the other. An irreducible quadratic Ax2 + 2Bxy + Cy2 + 2Dx + 2Ey + F = 0 is an equation of an ellipse provided that AC − B2 > 0. See ELLIPSE. Hyperbola. When > 1, the equation for the conic is written [x + c/( 2 − 1)]2 − y2/( 2 − 1) = 2c2/(1 − 2)2, ( > 1). Putting X = x + c/( 2 − 1), Y = y gives the standard form of a hyperbola X 2/A2 − Y 2/B2 = 1, where A = c/( 2 − 1) > 0, B = c/( 2 − 1)1/2 > 0. It is clear that the curve consists of two branches and is symmetric to the axes and the origin. Putting C = c 2( 2 − 1), then both F(C, 0), F(−C, 0) are foci of the hyperbola, with respective directrices X = c/( 2 − 1), X = −c/( 2 − 1), and A2 + B2 = C2. Then X 2/A2 − Y 2/B2 = 1 is seen to be an equation of the locus of points P(X, Y) such that PF − PF = 2A. The lines X/A ± Y/B = 0 are asymptotes of the hyperbola. As a point P traverses either branch of the hyperbola, its distances from an asymptote approach zero. An irreducible quadratic Ax2 + 2Bxy + Cy2 + 2Dx + 2Ey + F = 0 is an equation of an hyperbola provided AC − B2 < 0. See HYPERBOLA. Three-dimensional space. Let cartesian coordinate systems be established on each of three pairwise mutually perpendicular lines of three-space that intersect in O, the common origin of the systems. Suppose equal scales and call the lines the x axis, y axis, and z axis. To each point P of space an ordered triple (x, y, z) of real numbers is attached as rectangular cartesian coordinates, where x is the coordinate of the foot of the perpendicular from P to the x axis, and y and z are similarly defined. Thus every point of space has unique coordinates, and each ordered triple of real numbers is the coordinates of a point of space. If d denotes the distance of two points P1(x1, y1, z1), P2(x2, y2, z2), then d = [(x1 − x2)2 + (y1 − y2)2 + (z1 − z2)2]1/2. Let g denote any directed line and g the line through O parallel to g and directed in the same sense. If α, β, γ denote the angles that g makes with the x, y, and z axes, respectively (the direction angles of g), then cos α, cos β, cos γ are the direction cosines of g. They satisfy the relation cos2 α + cos2 β + cos2 γ = 1, and any three numbers λ, µ, ν such that λ2 + µ2 + ν 2 = 1 are the direction cosines of a (directed) line. Three numbers a, b, c proportional to the direction cosines λ, µ, ν of a directed line g are direction numbers of g. Clearly direction numbers of parallel lines are proportional. If P1(x1, y1, x1), P2(x2, y2, z2) are distinct points of g, then x2 − x1, y2 − y1, z2 − z1 are direction numbers of g. It follows that P(x, y, z) is on g if and only if
637
638
Analytic geometry x − x1 = t(x2 − x1), y − y1 = t(y2 − y1), z − z1 = t(z2 − z1), −∞ > t > ∞. These are parametric equations of g (t is the parameter), from which the symmetric equations (x − x1)/(x2 − x1) = (y − y1)/(y2 − y1) = (z − z1)/ (z2 − z1) follow at once. The direction cosines of g, directed from P1 to P2, are cos α = (x2 − x1)/d, cos β = (y2 − y1)/d, cos γ = (z2 − z1)/d, where d = [(x2 − x1)2 + (y2 − y1)2 + (z2 − z1)2]1/2. If α 1, β 1, γ 1, and α 2, β 2, γ 2 are direction angles of directed lines g1, g2, respectively, and θ denotes the angle between them, then cos θ = cos α 1 cos α 2 + cos β 1 cos β 2 + cos γ 1 cos γ 2. Hence g1, g2 are mutually perpendicular if and only if a1a2 + b1b2 + c1c2 = 0, where a1, b1, c1 and a2, b2, c2 are direction numbers of g1, g2, respectively. Let the plane π go through P0(x0, y0, z0) and be perpendicular to a line g with direction numbers a, b, c. Then P(x, y, z) is in π if and only if P = P0 or the line joining it to P0 is at right angles to a line through P0 that is parallel to g; that is, if and only if a(x − x0) + b(y − y0) + c(z − z0) = 0. Hence to each plane corresponds a linear equation. Conversely, if P0 (x0, y0, z0) satisfies the linear equation Ax + By + Cz + D = 0, with A, B, C not all zero, then A(x − x0) + B(y − y0) + C(z − z0) = 0 is an equation of a plane through P0, perpendicular to a line with direction numbers A, B, C, and so Ax + By + Cz + D = 0, with A, B, C = 0, 0, 0, is an equation of a plane. Clearly x = 0 is an equation for the plane determined by the y and z axes; equations of the other two coordinate planes are y = 0 and z = 0. If a directed perpendicular g from O to a plane meets the plane at P (g is directed from O to P in case the plane is not through O) the plane has equation x cos α + y cos β + zcos γ = p, where α, β, γ are the direction angles of g and p = OP. This is the normal form of equation of a plane. The general form is reduced to it upon dividing by 1 ±[A2 + B2 + C2] /2 . The distance from Ax + By + Cz + D + 0 to P(x0, y0) is (Ax0 + By0 + Cz0 + D)/±(A2 1 + B2 + C2) /2 , where the sign is selected opposite that of D. (In case D = 0, other conventions are used.) Two planes Aix + Biy + Ciz = 0, with i = 1, 2, are parallel in case the number triples A1, B1, C1 and A2, B2, C2 are proportional and are mutually perpendicular provided A1A2 + B1B2 + C1C2 = 0 (since A1, B1, C1 and A2, B2, C2 are direction numbers of lines that are perpendicular to the respective planes). If Pi(xi, yi, zi), with i = 1, 2, 3, is not collinear, the plane determined has the equation x y z 1 x y z 1 1 1 1 =0 x2 y2 z2 1 x y z 1 3
3
3
Three planes Aix + Biy + Ciz + Di = 0, with i = 1, 2, 3, intersect in one point if and only if A1 B2 C1 A2 B2 C2 = 0 A3 B3 C3 The locus of points whose coordinates satisfy an
equation f(x, y, z) = 0 is a surface. Curves may be thought of as intersections of two surfaces, and as such the coordinates of their points satisfy two equations f(x, y, z) = 0, g(x, y, z) = 0. Thus a line, considered as the intersection of two planes, is given by the two (simultaneous) equations Aix + Biy + Ciz = Di = 0, with i = 1, 2. That line has direction numbers C A A B B C 1 1 2 1 1 1 C2 A2 A2 B2 B2 C2 Curves are also represented parametrically by equations x = f(t), y = g(t), z = h(t), the parameter t varying in an interval (a, b), finite or infinite. Parametric equations for the line have already been given. Additional examples are x = r cos t, y = r sin t, z = 0, with 0 ≤ t ≤ 2π, parametric equations of the circle in the xy plane, with center at O and radius r; x = a · cos t, y = a · sin t, z = kt, with k = 0, − ∞ < t < ∞, parametric equations of a circular helix—the curve of the thread of a machine (untapered) screw. Special surfaces. It follows from the definition of a sphere and the formula for distance of two points that (x − a)2 + (y − b)2 + (z − c)2 = r2 is an equation for the sphere with center (a, b, c) and radius r, and by completing the squares of the x, y, and z terms in the equation x2 + y2 + z2 + 2Dx + 2Ey + 2Fz + G = 0, it is seen that the locus of such an equation is a sphere with positive or zero radius, or there is no (real) locus. Any equation in just two of the three coordinates is an equation of a cylinder whose elements are parallel to the axis of the missing variable. Thus the locus in 3-space of x2 + y2 = r2 is a (right) circular cylinder whose elements are parallel to the z axis and which intersects the xy plane in the circle x2 + y2 = r2, z = 0. Any equation f(x, y, z) = 0, with f(x, y, z) homogeneous in x, y, z (for example, 4xy − xz + yz = 0, x3 − xy2 + z3 = 0) has a cone with vertex O as locus. A surface of revolution is obtained by rotating a plane curve C about a line g of its plane. If f(x, y) = 0, z = 0 are equations of C, and g is the x axis, the resulting surface of revolution has the equation f(x, y2 + z2 ) = 0. Thus the surface generated by revolving the circle x2 + (y − b)2 = a2, z = 0, about the x axis (the torus or anchor ring, if b > a) has the equation x2 + ( y2 + z2 − b)2 = a2. A quadric surface is the locus of points whose coordinates satisfy an equation of the form Ax2 + By2 + Cz2 + Dxy + Exz + Fyz + Gx + Hy + Jz + K = 0, where at least one coefficient of a second-degree term is not zero. Some surfaces obtained by rotating conics about a line belong to this class, for example, spheres, prolate and oblate spheroids (given by rotating an ellipse about its major and minor axes, respectively), hyperboloids and paraboloids resulting from rotations of hyperbolas and parabolas about their axes of symmetry, and right circular cones and cylinders. Cylinders with conics for directrix curves are also members. Apart from degenerate cases, such as two planes, the remaining quadric surfaces and the standard forms of their equations are (1) ellipsoid,
Analytic hierarchy x2/a2 + y2/b2 + z2/c2 = 1; (2) hyperboloid of one sheet, x2/a2 + y2/b2 − z2/c2 = 1; (3) hyperboloid of two sheets, x2a2 − y2/b2 − z2/c2 = 1; (4) elliptic paraboloid, x2/a2 + y2/b2 = 2z; and (5) hyperbolic paraboloid, x2/a2 − y2/b2 = 2z. Hyperboloids of one sheet and hyperbolic paraboloids are ruled surfaces; that is, each contains an infinity of straight lines, called generators. In fact, each of those surfaces contains two sets of generators. See QUADRIC SURFACE. n-Dimensions. Let cartesian coordinate systems with equal scales be established on each of n pairwise mutually perpendicular lines intersecting in the common origin O, and label the lines OX1, OX2, . . . , OXn. To each point P of n space an ordered n-tuple (x1, x2, . . . , xn) of numbers is attached as coordinates, where xi is the coordinate of the foot of the perpendicular from P to OXi, with i = 1, 2, . . . , n. Two points P(x1, x2, . . . , xn), Q(y1, y2, . . . , yn) have distance d = [ (x1 − y2)2 + (x2 − y2)2 + · · · + (xn − yn)2]1/2. Direction angles x1, x2, . . . , xn of a directed line are defined as in three-dimensional analytic geometry, and the direction cosines satisfy the relation cos2 α 1 + cos2 α 2 + · · · + cos2 α n = 1. Direction cosines of the line g(P, Q), directed from P to Q, are (y1 − x1)/d, (y2 − x2)/d, . . . , (yn − xn)/d, and numbers proportional to them (for example, the numerators) are direction numbers of g(P, Q). Hence (X1, X2, . . . , Xn) are coordinates of a point on g(P, Q) if and only if (X1 − x1)/(y1 − x1) = (X2 − x2)/(y2 − x2) = · · · = (Xn − xn)/(yn − xn). These are symmetric equations for the line determined by P, Q. Denoting the common value of the quotients by t gives the n parametric equations Xi = xi + t(yi − xi), with i = 1, 2, . . . , n, −∞ < t< ∞. Let g be a line through C(ci, c2, . . . , cn) with direction numbers a1, a2, . . . , an. Point C, together with all points X such that the line g(C, X) is perpendicular to g, defines an (n − 1)dimensional subspace (hyperplane). Its equation is a1(X1 − c1) + a2(X2 − c2) + · · · + an(Xn −cn) = 0, since Xi − ci, with i = 1, 2, . . . , n, are direction numbers of g(C, X) and the equation is the condition that g and g(C, X) be mutually perpendicular. It is readily seen that the locus of every linear equation A1X1 + A2X2 + · · · + AnXn + K = 0 is a hyperplane perpendicular to a line with direction numbers A1, A2, . . . , An. It may be put in normal form by dividing by ±(A12 + A22 + · · · + An2)1/2, and the distance from it to a point C(c1, c2, . . . , cn) if found by substituting the coordinates of C for X1, X2, . . . , Xn. Subspaces of dimension k (1 ≤ k < n) are given by systems of n − k linear equations. An equation of the hyperplane determined by n points is readily expressed in determinant form, as well as the condition that n + 1 points be on a hyperplane [substitute the coordinates of the (n + 1)-st point in the first row of the determinant equation of the hyperplane]. Discussion of loci of higher order lies outside of the scope of this article. This brief sketch of analytic geometry has dealt only with that coordinate system most frequently used. For other coordinate systems (polar), as well as a discussion of transformation of coordinates, See AL-
GEBRA; COORDINATE SYSTEMS; DIFFERENTIAL GEOMETRY; EUCLIDEAN GEOMETRY; PROJECTIVE GEOMETRY; TRIGONOMETRY. Leonard M. Blumenthal
Bibliography. G. Fuller, Analytic Geometry, 7th ed., 1993; D. F. Riddle, Analytic Geometry, 6th ed., 1996.
Analytic hierarchy A framework for solving a problem. The analytic hierarchy process is a systematic procedure for representing the elements of any problem. It organizes the basic rationality by breaking down a problem into its smaller constituents and then calls for only simple pairwise comparison judgments, to develop priorities in each level. The analytic hierarchy process provides a comprehensive framework to cope with intuitive, rational, and irrational factors in making judgments at the same time. It is a method of integrating perceptions and purposes into an overall synthesis. The analytic hierarchy process does not require that judgments be consistent or even transitive. The degree of consistency (or inconsistency) of the judgment is revealed at the end of the analytic hierarchy process. Human reasoning. People generally provide subjective judgments that are based on feelings and intuition rather than on well-worked-out logical reasoning. Also, when they reason together, people tend to influence each other’s thinking. Individual judgments are altered slightly to accommodate the group’s logic and the group’s interests. However, people have very short memories, and if asked afterward to support the group judgments, they instinctively go back to their individual judgments. Repetition is needed to effect deep-rooted changes in people’s judgment. People also find it difficult to justify their judgments logically and to explicate how strong these judgments are. As a result, people make great compromises in their thinking to accommodate ideas and judgments. In groups, there is a willingness to compromise. If truth is to be an objective reality, reality must be very fuzzy because the search for truth often ends in compromise. What is regarded as truth may often be essentially a social product obtained through interaction rather than by pure deduction. Logical understanding does not seem to permeate judgment instantaneously. It apparently needs time to be assimilated. However, even when logical understanding has been assimilated, people still offer judgment in a spontaneous emotional way without elaborate explanation. Even when there is time, explanations tend to be fragmentary, disconnected, and mostly without an underlying clear logical foundation. Outline of the process. People making comparisons use their feelings and judgment. Both vary in intensity. To distinguish among different intensities, the scale of absolute numbers in Table 1 is useful. The analytic hierarchy process can be decomposed into the following steps. Particular steps may be emphasized more in some situations than in
639
640
Analytic hierarchy TABLE 1. Scale of relative importance Intensity of relative importance
Definition
Explanation
1
Equal importance
3 5
Slight importance of one over another Essential or strong importance
7
Demonstrated importance
9
Absolute importance
2, 4, 6, 8 Reciprocals of above nonzero numbers
Intermediate values between the two adjacent judgments If an activity has one of the above numbers assigned to it when compared with second activity, the second activity has the reciprocal value when compared to the first
others. Also as noted, interaction is generally useful for stimulation and for representing different points of view. 1. Define the problem and determine what knowledge is sought. 2. Structure the hierarchy from the top (the objectives from a broad perspective) through the intermediate levels (criteria on which subsequent levels depend) to the lowest level (which usually is a list of the alternatives). 3. Construct a set of pairwise comparison matrices for each of the lower levels, one matrix for each element in the level immediately above. An element in the higher level is said to be a governing element for those in the lower level since it contributes to it or affects it. In a complete simple hierarchy, every element in the lower level affects every element in the upper level. The elements in the lower level are then compared to each other, based on their effect on the governing element above. This yields a square matrix of judgments. The pairwise comparisons are done in terms of which element dominates the other. These judgments are then expressed as integers according to the judgment values in Table 1. If element A dominates element B, then the whole number integer is entered in row A, column B, and the reciprocal (fraction) is entered in row B, column A. Of course, if element B dominates element A, the reverse occurs. The whole number is then placed in the B,A position with the reciprocal automatically being assigned to the A,B position. If the elements being compared are equal, a one is assigned to both positions. The numbers used express an absolute rather than an ordinal relation. 4. There are n(n – 1)/2 judgments required to develop the set of matrices in step 3 (taking into account the fact that reciprocals are automatically assigned in each pairwise comparison), where n is the number of elements in the lower level. 5. Having collected all the pairwise comparison data and entered the reciprocals together with n unit entries down the main diagonal (an element is equal
Two activities contribute equally to the objective Experience and judgment slightly favor one activity over another Experience and judgment strongly favor one activity over another An activity is strongly favored and its dominance is demonstrated in practice The evidence favoring one activity over another is of the highest possible order of affirmation When compromise is needed
to itself, so a “one” is assigned to the diagonal positions), the eigenvalue problem Aw = λmaxw is solved and consistency is tested, using the departure of λmax from n (see below). 6. Steps 3, 4, and 5 are performed for all levels and clusters in the hierarchy. 7. Hierarchial composition is now used to weigh the eigenvectors by the weights of the criteria, and the sum is taken over all weighted eigenvector entries corresponding to those in the lower level of the hierarchy. 8. The consistency ratio of the entire hierarchy is found by multiplying each consistency index by the priority of the corresponding criterion and adding them together. The result is then divided by the same type of expression, using the random consistency index corresponding to the dimensions of each matrix weighted by the priorities as before. The consistency ratio should be about 10% or less to be acceptable. If not, the quality of the judgments should be improved, perhaps by revising the manner in which questions are asked in making the pairwise comparisons. If this should fail to improve consistency, it is likely that the problem should be more accurately structured; that is, similar elements should be grouped under more meaningful criteria. A return to step 2 would be required, although only the problematic parts of the hierarchy may need revision. If the exact answer in the form of hard numbers was actually available, it would be possible to normalize these numbers, form their ratios as described above, and solve the problem. This would result in getting the same numbers back, as should be expected. On the other hand, if firm numbers were not available, their ratios could be estimated to solve the problem. Example of the process. In the following example the analytic hierarchy process is used to assist a young family (a father, a mother, and a child) of specified income to buy a new car, say, either model A, B, or C. The choice will be determined through four important criteria.
Analytic hierarchy level 1: focus
TABLE 2. Comparison matrix comparing criteria for buying a car Decision to buy a car
Price
Price Running cost Comfort Status
1 1/3 1/7 1/8
Best New Car to Buy
Running cost Comfort Status Priorities 3 1 1/5 1/5
7 5 1 1/3
8 5 3 1
.582 .279 .090 .050 λ max = 4.198 C.I. = .066 C.R. = .073
The hierarchy of such a decision often takes the form shown in the illustration. In this hierarchy, level 1 is the single overall objective: Best New Car to Buy. On level 2 are the criteria which are perceived to compose what is meant by Best New Car, such as Price and Running Cost (operating and maintenance). On level 3 are the various alternative cars from which the family will choose. This downward decomposition format can easily be used on a wide class of problems. In addition, a slight further modification to incorporate feedback loops will cover an even wider range. The questions asked when comparing each criterion are of the following kind: Of the two alternatives being compared, which is considered more important by the family buying a car and how much more important is it? The comparison matrix of Table 2 is then formed. Since Price and Running Cost have the highest priorities, the other factors are discarded and only these two are used in continuing the process. Care must be taken in doing this. The original priorities are then normalized by dividing each of their sum to obtain the new relative priorities in Eqs. (1). .58 .58 + .28 .28 .33 = .58 + .28 .67 =
(1)
The process is then repeated for the third level, where each car is compared with respect to the two high-priority factors from the second level, as shown in Table 3. The two priority columns are then recorded as in Table 4. All entries of the first column are then multiplied by .67, the priority of Price, and those of the second column by .33, the priority of Running Cost, and added. This gives the third column. Thus Car B was selected for its efficient operation even
level 2: criteria
running cost
price
level 3: alternatives
comfort
Analytic hierarchy used to assist a family in buying a new car.
though its initial price is considerably higher than Car A. A car dealer in the income neighborhood of this family may find it profitable to stock up on the cars in the respective proportions. One of the most powerful contributions that the analytic hierarchy process makes is to test out the degree of inconsistency or incompatibility of new ideas or new policies adopted with older, more familiar, better tried successful methods. For example, in doing the above problem the participants were not sure whether the judgments for Price over Running Cost should be 7, 5, or 3. Each one was tried separately, and it was found that 3 yielded the highest consistency. Those who voted for 3 won that argument. Priorities and consistency. An easy way to get a good approximation of the priorities is to use the geometric mean. This is done by multiplying the elements in each row and taking their nth root, where n is the number of elements. Then, normalize the column of numbers thus obtained by dividing each entry by the sum of all entries. Alternatively, normalize the elements in each column of the matrix and then average each row. The consistency index can also be determined by hand calculations. Add the numbers in each column of the judgment matrix, multiply the first sum by the first priority, the second by the second, and so on, and add. For the first matrix the column sums (1.60, 4.40, 13.33, 17) are obtained, and multiplying by (.582, .279, .090, .050) gives 4.20. This number is
Car A Car B Car C
Car A
Car B
Car C
Priorities
Running cost
1 1/2 1/3
2 1 1/2
3 2 1
.540 .297 .163
Car A Car B Car C
λ max = 3.009 C.I. = .005 C.R. = .008
Car A 1 5 2
status
C
B
A
TABLE 3. Comparison matrices comparing alternative cars Price
641
Car B
Car C
Priorities
1/5 1 1/7
1/2 7 1
.106 .745 .150 λ max = 3.119 C.I. = .059 C.R. = .103
642
Analytic hierarchy TABLE 4. Composition of priorities
Car A Car B Car C
Price (priority .67)
Running cost (priority .33)
Composite priority of cars
.540 .297 .163
.106 .745 .150
.396 .445 .159
denoted by λmax. The consistency index is given by Eq. (2). C.I. =
λmax − n n−1
(2)
The consistency is now checked by taking the ratio (C.R.) of C.I. with the appropriate one of the set of numbers in Table 5 to see if it is about 10% or less (20% may be tolerated in some cases but not more). Otherwise the problem must be studied again and judgments revised. The consistency of the hierarchy in the above example, as given by Eq. (3), is .06, which is good.
ample, four identical chairs were placed at distances of 9, 15, 21, and 28 yards (1 yd = 0.9144 m) from a floodlight. The children stood by the light, looked at the line of chairs, and compared the first with the second, the first with the third and then with the fourth, and so on for the second, third, and fourth chairs. Each time the children said how much brighter one chair was, compared to the other. Their judgments were entered in the matrix of Table 6 to record the relative brightness of the chairs. The reciprocals were used in the transpose position. TABLE 6. Comparison matrix comparing perceived brightnesses of chairs
Chair 1 Chair 2 Chair 3 Chair 1 Chair 2 Chair 3 Chair 4
1 1/5 1/6 1/7
5 1 1/4 1/6
6 4 1 1/4
Brightness ratios
7 6 4 1
0.61 0.24 0.10 0.05 λ max = 4.39 C.I. = 0.13 C.R. = 0.14
.066 × 1 + .005 × .67 + .059 × .33 .089 = = .06 .900 × 1 + .580 × .67 + .580 × .33 1.48 (3) Judgment formation. When several people participate, judgments are often debated. Sometimes the group accepts a geometric average of their combined judgments. If there is strong disagreement, the different opinions can each be taken and used to obtain answers. Those which subsequently display the highest consistency within the group are the ones usually retained. The analytic hierarchy process incorporates equally both tangible factors, which require hard measurements, and such intangible factors as comfort, which require judgment. Eventually one finds that so-called hard numbers have no meaning in themselves apart from their utilitarian interpretations. In the above example, buying a $10,000 car is more than twice as “painful” as buying a $5000 car. The interdependence of criteria, such as Comfort and Price, has to be considered carefully since there may be some perceived overlap. For example, higher price buys more comfort, but it also buys other desirable attributes. Judging the relative importance of such things as price and comfort, therefore, must be done as independently as possible with avoidance of overlaps. Validation by physical laws. Using the scale 1–9 has been justified and demonstrated by many examples. However, the following simple optics illustration, carried out with small children, shows that perceptions, judgments, and these numbers lead to results which can be validated by laws of physics. In this ex-
Chair 4
The inverse-square law of optics is now used to test these judgments. Since the distances are 9, 15, 21, and 28 yd, these numbers are squared and their reciprocals calculated. This gives .0123, .0044, .0023, and .0013 respectively. Normalization of these values gives .61, .22, .11, and .06, which are very close to the brightness ratios obtained in the test using the analytic hierarchy process. Structuring a hierarchy. There are no rules for structuring a hierarchy. However, a typical analytic hierarchy for allocating resources—either by measuring costs or by measuring benefits—will often be stratified roughly as follows. The top level will include the overall objectives of the organization or system. Benefit-cost criteria may appear in the next level. A subordinate level may further clarify these criteria in the context of the particular problem by itemizing specific tasks which are to be accomplished at some level of performance. This is followed by the alternatives being evaluated. Capabilities. Designing an analytic hierarchy—like the structuring of a problem by any other method— is more art than science. It necessitates substantial knowledge of the system in question. A very strong aspect of the analytic hierarchy process is that the knowledgable individuals who supply judgments for the pairwise comparisons usually also play a prominent role in specifying the hierarchy. Although a hierarchy to be used in resource allocation will tend to have the vertical stratification
TABLE 5. Random consistencies n: Random consistency:
1 0
2 0
3 .58
4 .90
5 1.12
6 1.24
7 1.32
8 1.41
9 1.45
10 1.49
Analytical chemistry indicated above, it can also be much more general. The only restriction is that an element on a higher level must serve as a governing element for at least one element (which can be the element itself) on the immediately lower level. The hierarchy need not be complete; that is, an element at an upper level need not function as a criterion for all the elements in the lower level. It can be partitioned into nearly disjoint subhierarchies sharing only a common topmost element. Thus, for instance, the activities of separate divisions of an organization can be structured separately. As suggested above, the analyst can insert and delete levels and elements as necessary to clarify the task or to sharpen a focus on one or more areas of the system. The analytic hierarchy process has already been successfully applied in a variety of fields, including planning the allocation of energy to industries; designing a transport system for the Sudan; planning the future of a corporation and measuring the impact of environmental factors on its development; designing future scenarios for higher education in the United States; selecting candidates and winners in elections; setting priorities for the top scientific institute in a developing country; solving a faculty promotion and tenure problem; and predicting oil prices. See DECISION THEORY; SYSTEMS ENGINEERING. Thomas L. Saaty Bibliography. B. I. Golden (ed.), The Analytic Hierarchy Process, 1989; T. L. Saaty, The Analytic Hierarchy Process, 1980; T. L. Saaty, Decision Making for Leaders, 1982, reprint 1990; T. L. Saaty and L. Vargas, The Logic of Priorities, 1981, reprint 1991.
Analytical chemistry The science of chemical characterization and measurement. Qualitative analysis is concerned with the description of chemical composition in terms of elements, compounds, or structural units, whereas quantitative analysis is concerned with the measurement of amount. Scope. Originally, chemical analysis conformed closely to its literal meaning and consisted of the separation of a sample into its components, which were weighed. Later, such gravimetric techniques were supplemented by volumetric or, more appropriately, titrimetric analysis, which consists of measurement of the amount of a standardized solution of a reagent that reacts with a measured portion of the sample to reach an end point—a process known as titration. The end point is usually detected by means of an indicator, that is, a substance that changes color when a given amount of reagent has been added. These forms of chemical analysis are called classical or wet methods to distinguish them from the instrumental methods, which have greatly enlarged the scope of analytical chemistry. See GRAVIMETRIC ANALYSIS; TITRATION. Instrumental methods, which might better be termed physicochemical methods, have been developed to the degree that they make up the vast bulk
of both qualitative and quantitative analysis. They include measurements of purely physical characteristics, as well as the use of physical measurements, to follow the course of chemical reactions. Over the years, the scope of analytical chemistry has been enlarged from its original meaning, which was limited to the determination of chemical composition in terms of the relative amounts of elements or compounds in a sample. The discipline has been expanded to involve, for example, the spatial distribution of elements or compounds in a sample, the distinction between different crystalline forms of a given element or compound, the distinction between different chemical forms (such as the oxidation state of an element), the distinction between a component on the surface or in the interior of a particle, and the detection of single atoms or molecules of a substance. To permit these more detailed questions to be answered, as well as to improve the speed, accuracy, sensitivity, and selectivity of traditional analysis, a large variety of physical measurements are used. These methods are based on spectroscopic, electrochemical, chromatographic, chemical, and nuclear principles. Modern analysis has also placed significant demands on sampling techniques. It has become necessary, for example, to handle very small liquid samples [in the nanoliter (10−9 liter) range or less] as part of the analysis of complex mixtures such as biological fluids, and to simultaneously determine many different components. The sample may be a solid that must be converted through vaporization into a form suitable for analysis. Spectroscopy. One important area is spectroscopy, which includes the measurement of emission, absorption, reflection, and scattering phenomena resulting from interaction of a sample with gamma rays and x-rays at the high-energy end of the spectrum and with the less energetic ultraviolet, visible, infrared, and microwave radiation. See SPECTROSCOPY. Molecular spectroscopy. Lower-energy forms of excitation such as ultraviolet, visible, or infrared radiation are used in molecular spectroscopy. Ultraviolet radiation and visible radiation, which are reflective of the electronic structure of molecules, are used extensively for quantitative analysis. The radiation absorbed by the sample is measured. It is also possible to measure the radiation emitted (emission, fluorescence). The absorption of infrared radiation is controlled by the properties of bonds between atoms, and accordingly it is most widely used for structure identification and determination. It is less widely used for quantitative analysis except for gases such as carbon monoxide (CO) and hydrocarbons. X-rays are used through emission of characteristic radiation, absorption, or diffraction. In the last case, characteristic diffraction patterns reveal information about specific structural entities, such as a particular crystalline form. Extended x-ray absorption fine structure (EXAFS) is based on the use of x-rays from a synchrotron source to reveal structural details such as interatomic distances. See EMISSION SPECTROCHEMICAL ANALYSIS; EXTENDED X-RAY ABSORPTION
643
644
Analytical chemistry FINE STRUCTURE (EXAFS); INFRARED SPECTROSCOPY; X-RAY FLUORESCENCE ANALYSIS. Mass spectrometry. Mass spectrometry is an impor-
tant and increasingly applied method of analysis, especially for organic and biological samples. Among the applications are the analysis of more than 70 elements (spark-source mass spectrometry), surface analysis (secondary ion mass spectrometry and ionprobe mass spectrometry), and the determination of the structure of organic molecules and of proteins and peptides (high-resolution mass spectrometry, Fourier-transform mass spectrometry). See MASS SPECTROMETRY; SECONDARY ION MASS SPECTROMETRY (SIMS). Nuclear magnetic resonance. This technique measures the magnetic environment around individual atoms and provides one of the most important means for deducing the structure of a molecule. Atoms possessing nuclear spin are probed by monitoring the interaction between their nuclear spin and an applied external magnetic field. For large molecules, these interactions are complex, and a variety of nuclear excitation techniques have been developed that permit establishment of the connectivity between the various atoms in a molecule. Since the technique is nondestructive, it can be used to monitor living systems. See NUCLEAR MAGNETIC RESONANCE (NMR). Electron spectroscopy. Several forms of spectroscopy are especially useful for surface analysis. The scanning electron microscope (SEM) involves a finely collimated electron beam that sweeps across the surface of a sample to produce an image. At the same time, the surface atoms are excited to emit characteristic x-rays, making it possible to obtain an image of the surface along with its spatially resolved elemental composition. The resolution of this technique (electron microprobe) is in the micrometer (10−4 cm) range. Images with a resolution of nanometers (10−9 cm) have been obtained by using the techniques of atomic force microscopy (AFM) and scanning tunneling microscopy (STM), and this resolution corresponds to the dimensions of individual atoms. A significant advantage of the latter two techniques is that a high vacuum is not required, so samples can be analyzed at atmospheric pressure. See ELECTRON-PROBE MICROANALYSIS; ELECTRON SPECTROSCOPY; SCANNING ELECTRON MICROSCOPE; SCANNING TUNNELING MICROSCOPE. There are a number of other important surface analysis techniques involving the interaction of radiation with surface atoms. Low-energy electron diffraction (LEED) involves the observation of the diffraction pattern of a reflected electron beam to give structural information about the crystallographic orientation of the atoms on the surface. Several techniques are used to determine composition and structural information by measuring the energies of electrons emitted from surface atoms. These include ultraviolet photoelectron spectroscopy (UPS) and x-ray photoelectron spectroscopy (also called electron spectroscopy for chemical analysis, or ESCA), which differ in the energy of the photon beam striking the surface. Other techniques include
Auger electron spectroscopy (AES) and electron energy loss spectroscopy (EELS). See AUGER EFFECT; ELECTRON DIFFRACTION; SURFACE AND INTERFACIAL CHEMISTRY. Radiochemistry. The radiation emitted by a radionuclide can be determined accurately, forming the basis of a number of important analytical techniques. In some cases, these nuclides are incorporated directly into a molecule, rendering it radioactive [carbon-14 (14C), hydrogen-3 (3H), iodine-125 (125I), phosphorus-32 (32P)]. In other cases, the radioactivity is induced in the sample by bombardment with neutrons or photons (neutron or photon activation analysis). The radionuclides are identified by their characteristic radiation. For example, gammaray spectrometry is used to measure the characteristic radiation from nuclides emitting gamma rays. Many nuclides can be analyzed simultaneously by this method. See ACTIVATION ANALYSIS; GAMMA-RAY DETECTORS; RADIOACTIVITY. Electrochemical analysis. There are a number of widely applied techniques that are based mainly on four electrochemical phenomena: (1) potentiometry, involving the relationship between the measurement of an electrical potential between two electrodes, one of which responds selectively to the concentration (activity) of a species in solution; (2) amperometry, involving the measurement of a current proportional to the concentration of a species in solution at a given applied potential between two electrodes; (3) coulometry, based on the relationship between the charge passed in an electrochemical cell and the amount of a species generated (Faraday’s law); and (4) conductimetry, involving the relationship between the impedance (resistance) of a solution and the concentration of ionic species present. See ELECTROCHEMISTRY. Potentiometry. This is undoubtedly the most widely applied electrochemical technique, since it includes a variety of ion-selective electrodes, the most important of which is the glass electrode used to measure pH. Other important ion-selective electrodes measure ions of sodium, potassium, calcium, sulfide, chloride, and fluoride. When the electrodes are used in conjunction with gas-permeable membranes, gases such as ammonia, carbon dioxide, and hydrogen sulfide can be measured. See ION-SELECTIVE MEMBRANES AND ELECTRODES; pH; POLAROGRAPHIC ANALYSIS. Amperometry. This method depends on the detection of species that can undergo electrolysis; many substances such as oxygen are detected in this way. If the detection electrode is covered with an enzyme layer, a variety of substances in complex mixtures can be determined by taking advantage of the selectivity of the enzymatic reaction. An electroactive product of the enzymatic reaction is detected, and this configuration is called an enzyme electrode. These determinations are typically performed at a constant applied potential. A number of important techniques are based on measuring the current resulting from time-dependent variation of the potential. These techniques include cyclic
Analytical chemistry voltammetry, chronocoulometry, pulse techniques (such as pulsed amperometric detection), and chronoamperometry. They are used to study electrochemical reactions at electrodes, to study properties of electrode surfaces, and to analyze (for example, to detect an effluent from a chromatographic column or to determine a trace metal in solution). Other phenomena. Coulometry is used in quantitatively generating analytical reagents and in connection with a system capable of detecting an end point. Conductimetry is a nonselective technique that is used to measure the overall ionic content of an electrically conducting solution. See COULOMETER. Herbert A. Laitinen; George S. Wilson Trace analysis. Trace analysis involves the determination of those elemental constituents of a sample that make up approximately 0.01% by weight of the sample or less. There is no sharp boundary between nontrace and trace constituents. The lower limit of detected concentration is set by the sensitivity of the available analytical methods and, in general, is pushed downward with progress in analytical techniques. A large number of different physical and chemical techniques have been developed for the measurement of the elemental composition at the microgram-per-gram and nanogram-per-gram level, thereby constituting the field of trace analysis. Methods. All trace analytical methods can be divided into three component steps: sampling, chemical or physical pretreatment, and measurement. Depending upon the type of information desired in an analysis and the requirements of sensitivity, precision, and other performance figures of merit, an appropriate measurement technique is selected. Most frequently, this will be an instrumental approach. Therefore, it is essential to know the capabilities and limitations of the various methods of instrumental measurement available for trace determinations. Once they are known, appropriate steps can be taken in the sampling and pretreatment steps to provide sufficient amounts of the microconstituents that are free of interferences and in the appropriate form for the final measurement. In a number of methods the pretreatment step may be omitted, and in others the sampling and measurement occur simultaneously. In spite of possible deviations, these steps are interrelated and require different degrees of emphasis, depending upon the individual analytical situation. In many cases, the analyst uses a particular physical or physicochemical method in which manifestations of energy provide the basis of measurement. These methods are indirect in the sense that the emission or absorption of radiation or transformation of energy must be related in some way to the mass or concentration of the species that are being determined. The establishment of these relations almost invariably requires calibration, with the use of standards of known content of the constituent in question. As such, many of the available techniques do not provide absolute results. Trace analysis ranges from the more classical chemical methods of colorimetric and absorption spectrophotometric analysis to modern instrumen-
Some trace analysis methods and limits of detection
Method Chromatography Thin-layer Gas-liquid Liquid-liquid Electrochemical Coulometry Ion selective electrode Polarography Conventional Modern
Absolute, g
Limits of detection: concentration, ppm
10−5 to 10−3 10 to 106 10−3 to 100 10−9 to 100
10−2 to 102 100 to 103 10−3 to 103
Laser probe microanalysis
102 to 104
Nuclear Neutron activation Electron probe Ion probe
10−3 to 10−1 102 to 103 10−1 to 101
Mass spectrometry Isotope dilution Spark source
10−5 to 106 10−3 to 101 10−5 to 100
Electrical plasma* Organic microanalysis Optical-absorption Atomic Flame Nonflame Molecular UV-visible Infrared Microwave
>10
−5
10−15 to 10−9
10−3 to 102 103 to 106 100 to 103
Optical-emission AC spark DC arc Electrical plasma† Optical-fluorescence Atomic Flame Nonflame Molecular
10−3 to 101
101 to 103 10−2 to 102 10−4 to 102
10−15 to 10−9
10−3 to 102 10−3 to 10−1
Optical-phosphorescence
10−3 to 102
Optical-Raman
100 to 105
Spectrometric-resonance Nuclear magnetic Electron spin
10−9 to 10−6
Thermal analysis
10−5 to 10−4
Wet chemistry Gravimetry Titrimetry
10−3 to 10−2
X-ray spectrometry Auger ESCA Fluorescence Mossbauer ¨ Photoelectron Immunoassay
101 to 105
10−2 to 104 103 to 105 103 to 105 10−1 to 102 100 to 103 100 to 103
10−2 to 10−9
10−3 to 100
*Inductively coupled plasma or microwave-induced plasma. †Inductively coupled plasma or direct-current plasma.
tal approaches. The wide diversity of methods is apparent from the number of approaches and the attendant classes and subclasses in the table. Within any one of these categories, there are techniques that are more specialized. Of the various criteria used in the selection of an
645
646
Analytical chemistry appropriate trace analytical method, sensitivity, accuracy, precision, and selectivity are of prime importance. Other important considerations, such as scope, sampling and standards requirements, cost of equipment, and time of analyses, are of great practical significance. Sensitivity and detection limits. For all analysis regimes, the measure x of some physical parameter is related to the concentration c of the analyte in a certain sample. The sensitivity S is the slope of the analytical calibration curve, which is the plot of the measure x versus the concentration of the analyte in a series of standards having known analyte concentrations. The term sensitivity must not be used to indicate either the limit of detection or the concentration required to give a certain signal. The standard deviation σ is given by Eq. (1), where n
(xi − x¯ )2 σ = (1) n−1 i=1 xi is the individual measure, and x¯ is the mean of n measurements. If n is less than about 10, σ is replaced by the term s. The relative standard deviation (RSD) is given by Eqs. (2). RSD =
σ x¯
or
RSD =
s x¯
(2)
Precision is defined as the random uncertainty for the measure x of the corresponding uncertainty in the estimate of the concentration. It is often expressed as the relative standard deviation. Accuracy is a measure of the agreement between the estimated concentration and the true value. Accuracy can never be better than the precision of the analytical procedure. Bias (systematic errors) in the calibration procedure always causes the estimated uncertainties to disagree with the true value by an amount equal to the bias. The limit of detection cL is the smallest concentration that can be detected with a reasonable certainty. This concentration is given by Eq. (3), where cL =
xL − x¯b S
(3)
xL is the limiting detectable measure and xb is the average blank measure. The value for xL is given by Eq. (4), where k is a factor determining a margin of xL = xb + ksb
(4)
confidence and sb is the standard deviation of a limited number of measures of the blank. Generally, k is taken as 3; thus, the value of cL can be determined by Eqs. (5). This means that any concentration recL =
3sb S
or
cL =
3Nb S
(5)
sulting in a signal three times the background standard deviation is considered just detectable. To evaluate xb and sb, at least 20 measurements should be used. If the major sources of variation are electrical
noises, sb can be replaced with Nb, the background noise level. If 3sb is chosen, the confidence level is 99.86% for a purely Gaussian distribution of errors. At low concentrations, broader and asymmetric distributions are likely, so 3sb corresponds to a practical confidence level of about 90%. It is noted that at the limit of detection the relative standard deviation is about 0.5; that is, there is equal confidence in deciding whether the analyte is detected or not. As a result of widely varying pathways of development of many analytical techniques, there is a lack of consistency in the definition or specification of detection limits in the analytical literature. Consequently, it is difficult to make critical comparisons between methods for a particular element. However, since a particular element may be determined by a number of different techniques, depending upon the matrix in which it is being sought, a summary of the experimental values that have been published is pertinent. See ACTIVATION ANALYSIS; ANALYTICAL CHEMISTRY; CHROMATOGRAPHY; ELECTROCHEMICAL TECHNIQUES; GAS CHROMATOGRAPHY; SPECTROSCOPY. Andrew T. Zander Separation techniques. This collection of techniques includes various forms of chromatography and electrophoresis. They are based on the separation of a mixture of species in a sample due to differential migration. Two forces act in opposition: a stationary phase acts to retard a migrating species, while the mobile phase tends to promote migration. The mobile phase may be liquid (liquid chromatography) or gaseous (gas chromatography), while the stationary phase may be a solid or a solid covered with a thin film of liquid. The stationary phase is typically packed in a column through which the mobile phase is pumped. High-performance liquid chromatography (HPLC) has become especially important for the separation of complex mixtures of nonvolatile materials, because the use of high pressure emphasizes the differences in the rate of diffusion among the species to be separated. Separations may often be accomplished in a matter of several minutes. The stationary phase can preferentially interact with the migrating species according to charge, size, and hydrophobicity, or in some cases because of the special affinity that a species has for the stationary phase (affinity chromatography). The stationary phase can also be a thin layer of solid support deposited on a plate (thin-layer chromatography). See CHROMATOGRAPHY; GAS CHROMATOGRAPHY; LIQUID CHROMATOGRAPHY. Alternatively, the driving force for separation will be the migration of charged species in an electric field (electrophoresis). The stationary phase may be a gel on a plate or in a tube, or a solution maintained in a capillary through which the analytes move. The important techniques in this area are capillary electrophoresis, isotachophoresis, and isoelectric focusing. See ELECTROPHORESIS. In order to detect the various species after they are separated, most of the techniques described above have been employed. In some cases, it is desir-
Anamnia able to have detection methods that are specific for the species of interest. In other cases, such specificity is not necessary since the species have already been separated and a detector with a general response can be used. Chemical methods. Chemical and biological reactions form the basis of many analytical methods, besides the classical forms of gravimetric and volumetric analysis. From biology have come a number of widely used techniques: immunoassays, enzymecatalyzed reactions, gene probes, and bioluminescence. These methods are often based on the rate of a reaction, which is proportional to the concentration of the analyte of interest. See BIOLUMINESCENCE; IMMUNOASSAY; KINETIC METHODS OF ANALYSIS. Thermal methods. Thermal methods are based on the heating of a sample over a range of temperatures. This approach may result in absorption of heat by the sample or in evolution of heat due to physical or chemical changes. Thermogravimetry involves the measurement of mass; differential thermal analysis involves a detection of chemical or physical processes through a measurement of the difference in temperature between a sample and a stable reference material; differential thermal calorimetry evaluates the heat evolved in such processes. A variety of calorimetric techniques are used to measure the extent of reactions that are otherwise difficult to evaluate. See CALORIMETRY. Sample handling and processing methods. As the demands increase for more rapid and reliable analysis at lower cost, automation of sampling and sample handling, measurement, and data analysis have become extremely important. Applications may involve monitoring a manufacturing process stream, air quality, or critical diagnostic analyses of a patient undergoing surgery. Fortunately, the advent of low-cost but powerful computers has greatly aided development in this area. Computers. The digital computer has had a major impact on modern analytical chemistry. The microprocessor has been incorporated as an integral part of a variety of analytical instruments to permit detailed programming of their operation. More complex programming and automation are possible through the personal computer, which is being used not only to operate a single instrument but also to tie various operations together. These systems are increasingly linked to networks, which provide for storage of and access to data. See DIGITAL COMPUTER; MICROCOMPUTER; MICROPROCESSOR. Chemometrics. Chemometrics refers to the use of statistical and other mathematical approaches to data handling. An important application is the optimization of an instrumental procedure by processing the initial instrument output rapidly and feeding the information back in order to control the settings of the instrument, thereby achieving optimal accuracy and sensitivity from the instrument. Still more complex computer operations, such as pattern recognition, are designed to draw conclusions through comparisons of arrays of data. See CHEMOMETRICS.
Automation. Analytical methods can be automated by several approaches. Clinical methods are often carried out by means of an automated analyzer. The sample is introduced into a flowing stream (unsegmented flow) or separated from the next sample by an air bubble (segmented flow), and a sequence of analytical operations is performed on the sample, culminating in a final measurement step based on an instrumental reading. The instrument is calibrated by running a standard sample through the system. A computer readout compares the analytical result with the normal range expected for the analysis. Other physical principles are also used in automated analytical systems. In robotics, a computercontrolled robot carries out the routine steps of chemical methods, including sample preparation and cleanup, weighing, dissolution, adjustment of conditions, reagent manipulation, and instrument reading. This capability increases sample throughput and reduces errors caused in many cases by operator boredom and inattention due to the repetitive nature of the manipulations. Such analyses can be performed very rapidly with small volumes (nanoliters) of solution. See COMBINATORIAL CHEMISTRY; ROBOTICS. Instrumentation. As there is a constant need for improved instrumentation, the development of new instrumentation, rather than simple application of old techniques, is an area of active interest in analytical chemistry research. See CHEMISTRY; INSTRUMENT George S. Wilson SCIENCE. Bibliography. D. C. Harris, Quantitative Chemical Analysis, 6th ed., 2002; A. G. Howard and P. J. Statham, Inorganic Trace Analysis: Philosophy and Practice, 1994; J. R. Lawrence (ed.), Trace Analysis, vol. 1, 1981, vol. 2, 1982; J. F. Rubinson and K. A. Rubinson, Contemporary Chemical Analysis, 1998; D. A. Skoog et al., Principles of Instrumental Analysis, 5th ed., 1998.
Anamnia Those vertebrate animals, sometimes called Anamniota, which lack an amnion in development. The amnion is a protective embryonic envelope that encloses the embryo and its surrounding liquid, the amniotic fluid, during fetal life. An amnion is present in mammals, birds, and reptiles, but is absent in fishes and amphibians. In early classifications the vertebrates were commonly separated on this basis, and the expressions Amniota and Anamniota are still useful in grouping higher and lower vertebrates. It should be recognized that these terms represent grades of development, however, and do not carry the connotation of established classificatory ranks. Anamnia, then, is a group name that includes the Recent classes Agnatha, Chondrichthyes, Osteichthyes, and Amphibia and, by presumption, the class Placodermi, which is known only from fossils. See AMNION; AMNIOTA; AMPHIBIA; PISCES (ZOOLOGY). Reeve M. Bailey
647
648
Anaphylaxis
Anaphylaxis A generalized or localized tissue reaction occurring within minutes of an antigen-antibody reaction. Similar reactions elicited by nonimmunologic mechanisms are termed anaphylactoid reactions. The term anaphylaxis was introduced by C. R. Richet to describe acute adverse reactions which followed reinjection of an eel toxin into dogs. Instead of becoming immune (nonreacting) to the toxin, the animals developed severe and often fatal congestion and hemorrhage of the gastrointestinal tract, lungs, pleura, and endocardium. It was initially thought that anaphylaxis in humans was associated with precipitating antibody (known today to be of the IgG class of immunoglobulins), and that it differed from the common allergic manifestations of hay fever due to “reaginic” antibody. The reaginic antibody has now been identified as IgE. See ALLERGY; IMMUNOGLOBULIN. Signs and symptoms. In animals undergoing anaphylaxis, one organ system is usually affected more than others (thus the term shock-organ). In the guinea pig, the manifestations of anaphylaxis are respiratory, with cough, wheezing, and gasping respirations caused by severe bronchoconstriction. In the dog and the rat, the manifestations are gastrointestinal with hemorrhage. Severe pulmonary vasoconstriction leading to heart failure and death prevail in the rabbit. In humans, the clinical manifestations of anaphylaxis include reactions of the skin with itching, erythema, and urticaria; the upper respiratory tract with edema of the larynx; the lower respiratory tract with dyspnea, wheezing, and cough; the gastrointestinal tract with abdominal cramps, nausea, vomiting, and diarrhea; and the cardiovascular system with hypotension and shock. Individuals undergoing anaphylactic reactions may develop any one, a combination, or all of the signs and symptoms. Anaphylaxis may be fatal within minutes, or may occur days or weeks after the reaction, if the organs sustained considerable damage during the hypotensive phase. Tissue changes in acute anaphylaxis are minimal and include edema and a mild eosinophilic infiltration. Myocardial ischemia and ischemic changes of the renal tubules and the brain may be seen after prolonged hypotension. Macroscopic changes are variable, and in fatal cases usually involve the respiratory system, where pulmonary edema, pulmonary infiltrates, and increased bronchial secretions are noted. Mechanism. Anaphylaxis in humans is most often the result of the interaction of specific IgE antibody fixed to mast cells and antigen. Two molecules of IgE are bridged by the antigen, which may be a complex protein or chemical (hapten) bound to protein. The antigen-antibody interaction leads to increased cellmembrane permeability, with influx of calcium and release of either preformed or newly formed pharmacologic mediators from the granules. Preformed mediators include histamine and eosinophilic or neutrophilic chemotactic factors. Newly formed molecules include leukotrienes or slow-reacting
substance of anaphylaxis and prostaglandins. The mediator action induces bronchoconstriction, vasodilation, cellular infiltration, and increased mucus production. Degranulation of mast cells and basophils is augmented when intracellular cyclic guanine monophosphate (cGMP) is increased; the reaction is inhibited in the presence of increasing intracellular concentrations of cyclic adenine monophosphate (cAMP). See EICOSANOIDS. Another mechanism for induction of anaphylaxis in humans occurs when antigen binds to preformed IgG antibody and complement components interact with the antigen-antibody complex. The early components of the complement system bind to the antibody molecule, leading to activation of other complement components. During the activation, components known as anaphylatoxins (C3a and C5a) are released which may directly cause bronchoconstriction with respiratory impairment, and vasodilation with hypotension or shock. Activation of the complement system by agents like iodinated radiocontrast media, or degranulation of mast cells by drugs like opiates, will elicit similar manifestations without an antibody-antigen interaction. See COMPLEMENT. Anaphylaxis due to IgE mechanisms has been associated with foreign proteins such as horse antitoxins, insulin, adrenocorticotropic hormone (ACTH), protamine, and chymopapain injected into herniated discs; drugs such as penicillin and its derivatives; foods such as shellfish, nuts, and eggs; and venom of stinging insects. Anaphylaxis mediated by IgG is seen in blood-transfusion reactions and following the use of cryoprecipitate, plasma, or immunoglobulin therapy. Direct mast-cell degranulation with release of mediators occurs with opiates, radiocontrast media, and curare. In certain individuals, mast cell degranulation occurs after exercise and following exposure to cold. Individuals with systemic mastocytosis spontaneously release mediators from the increased number of tissue mast cells, which may cause an anaphylactic reaction. When no cause for the anaphylaxis can be found, the reaction is termed idiopathic. Prevention and treatment. After the identification of the inciting agent for the anaphylactic reaction, prevention is the best mode of therapy. Individuals should not be given antibiotics to which they are sensitive, and foods which cause reactions should be avoided. Individuals sensitive to insect-sting venom should take preventive measures to avoid stings and should use insect repellents. Individuals who react to exercise should avoid such activities; when cold is the cause of anaphylaxis, warm clothes should be worn, and submerging in the cold water of swimming pools should be avoided. Immunotherapy with insect venom and desensitization with certain drugs are effective prophylactic measures. Individuals with recurrent episodes of anaphylaxis, when the etiological cause is unknown and preventive measures are impractical, should be provided with epinephrine in a form that can be self-administered whenever symptoms occur.
Anaplasmosis Administration of corticosteroids may also be helpful. Such prophylaxis may be life-saving. See EPINEPHRINE. The treatment of anaphylaxis is aimed at reducing the effect of the chemical mediators on the end organs and preventing further mediator release. The drug of choice for this is epinephrine given subcutaneously in repeated doses. Additionally, a clear airway and appropriate oxygenation must be maintained; hypotension should be treated, as should any cardiac arrhythmia. Aminophyllin, a rapid-onset bronchodilator, and corticosteroids, which alter prostaglandin and leukotriene metabolism but have an onset of action in several hours, are second-line indications to be used in patients who do not respond adequately to epinephrine. Cromolyn sodium, known to prevent mast-cell degranulation, has been effective in preventing bronchoconstriction, but has no place in the treatment of the acute anaphylactic reaction. Such effective therapeutic measures must be carried out as soon as possible because of the potentially fatal nature of anaphylaxis. See HYPERSENSITIVITY; SHOCK SYNDROME. Shlomo Bar-Sela; Jordan N. Fink Bibliography. E. Bacal, R. Patterson, and C. R. Zeiss, Evaluation of severe (anaphylactic) reactions, Clin. Allergy, 8:295, 1978; E. Middleton, C. E. Reed, and E. F. Ellis (eds.), Allergy: Principles and Practices, 4th ed., 1993; R. Patterson (ed.), Allergic Diseases: Diagnosis and Management, 5th ed., 1997; P. L. Smith et al., Physiologic manifestations of human anaphylaxis, J. Clin. Invest., 66:1072–1080, 1980.
Anaplasmosis A disease of mammals caused by a specialized group of gram-negative bacteria in the order Rickettsiales, family Anaplasmataceae, genus Anaplasma. Anaplasma is an obligate intracellular parasite that infects erythrocytes (red blood cells) of cattle, sheep, goats, and wild ruminants in much of the tropical and subtropical world. A newly reclassified species (Anaplasma phagocytophilum) that infects the white blood cells of mammals was recently added to this genus. However, the most important species and the one traditionally associated with the disease anaplasmosis is A. marginale. Anaplasmosis is one of the most important diseases of cattle, causing anemia and sometimes death, resulting in significant economic losses worldwide. Occurrence and transmission. Anaplasma resides in erythrocytes within a vacuole formed from the host cell membrane. The rickettsiae remain within the vacuole and multiply by binary fission, producing numerous rickettsiae. The vacuole—which can be identified in stained blood films as a round, densely staining body approximately 1 micrometer in diameter—is referred to as an inclusion body. Anaplasma marginale contains both deoxyribonucleic acid (DNA) and ribonucleic acid (RNA), with a total genome size of 1200–1260 kilobase pairs (kbp). Molecular studies based on the analysis of
ribosomal RNA relate Anaplasma closely to Ehrlichia spp. Transmission of Anaplasma occurs biologically by multiple species of ticks. In the United States Dermacentor spp ticks are the primary vectors, but some 20 species act as vectors worldwide. Rickettsiae are ingested with a blood meal and undergo complex development beginning in the tick gut cells. Anaplasma is transmitted from the salivary glands while ticks feed on the host. Infections acquired by one tick stage may be transmitted after molting to the next stage (transtadial). Ticks may remain infected indefinitely (without feeding) and serve as a reservoir for infection of susceptible cattle. Mechanical transmission occurs when the mouthparts of biting flies (such as horse flies and stable flies) become contaminated while feeding on infected cattle and then quickly move to uninfected animals. There is no development of rickettsiae in these mechanical vectors. Contaminated needles and instruments for dehorning, castration, and tagging may also transmit the organism throughout a herd. Isolates of Anaplasma from different geographic areas differ from each other in biology, morphology, and antigenic characteristics. Epidemiology. Once cattle are infected with Anaplasma, they remain persistently infected for life and may serve as reservoirs of infection for other cattle or tick vectors. During this time the parasitemia (presence of parasites in the blood) fluctuates and at times may be undetectable. The spread of anaplasmosis may occur when these carrier cattle are moved to nonendemic areas. Ticks have been shown to transmit new infections after feeding on animals with undetectable parasitemias. Carrier animals may also have relapse infections, producing new infections in susceptible herds. Five major surface proteins of Anaplasma have been identified, and some have been found to be encoded by polymorphic multigene families. This variability results in distinct geographic strains of Anaplasma which are antigenically different and complicate vaccine strategies. Clinical disease. Once an animal is infected it takes approximately 21 days before organisms can be identified in blood smears; but animals may not exhibit clinical disease for as long as 60 days postinfection. In acute anaplasmosis, cattle exhibit depression, loss of appetite, increased temperature, labored breathing, dehydration, jaundice, and a decrease in milk production. Erythrocyte count, packed cell volume, and hemoglobin values decrease significantly, and death may result from severe anemia. Abortions may also occur following infection. These clinical manifestations of anemia are associated with the rapid removal of infected and uninfected erythrocytes by macrophages (phagocytic cells of the immune system). Recovered animals gradually regain condition but remain chronically infected and are subject to periodic relapses. They may also serve as a source of infection for susceptible animals. The disease is most severe in older animals not previously exposed; calves up to 1 year old become infected but usually
649
650
Anapsida have a mild or subclinical (showing no symptoms) reaction. Diagnosis. Parasites can be identified as a round inclusion body at the periphery of red blood cells in a Giemsa-stained blood film. Infection can be confirmed by identifying specific antibodies in serological tests such as the competitive enzyme-linked immunosorbent assay (C-ELISA). Molecular amplification techniques such as the polymerase chain reaction (PCR) are also used, particularly when parasite density is very low or undetectable. Immunity. Cattle that recover from clinical anaplasmosis are protected from subsequent exposure to the same geographic isolate of Anaplasma. Calves have a natural immunity to clinical disease if exposed within their first year. Protection is attributed to a complex combination of antibody and cell-mediated responses, with initial responses directed against the outer-membrane proteins of Anaplasma. Animals develop high levels of antibody during acute and recovery phases of anaplasmosis, but passive transfer of the antiserum (the serum component of blood that contains antibodies specific to one or more antigens) does not convey protection to naive animals (which have not been previously exposed to the disease). Splenectomy of recovered animals results in recrudescence (return of symptoms) of the rickettsemia (presence of rickettsiae in the blood), suggesting a role for cellular immunity in protection. Cell-mediated mechanisms of protection may involve macrophage activation and specific T-lymphocyte responsiveness. Specific antibodies may coat infected erythrocytes, enhancing their susceptibility to phagocytosis by activated macrophages. The ability of the rickettsiae to persist within the host at low levels despite a strong immune response is thought to be the result of antigenic variants created during cyclical rickettsemia. See ANTIBODY; ANTIGEN; CELLULAR IMMUNOLOGY; IMMUNITY. Control and treatment. The clinical manifestations of anaplasmosis can be halted or prevented if tetracycline antibiotics are administered early. The drug inhibits rickettsial protein synthesis but does not kill the organism. Tetracyclines have been added to feeds to prevent clinical symptoms, but this does not prevent infection. Controlling fly and tick populations by chemical spraying or dipping have also been used to reduce the spread of disease. Vaccination with killed A. marginale from erythrocytes has been shown to provide some protection against the acute disease but it does not prevent infection. Blood-derived vaccines are no longer used in most of the United States. The related and less virulent species, A. centrale, is used as a liveblood vaccine in some countries and provides some protection against A. marginale, but it can revert to a virulent form. Recombinant surface proteins of A. marginale have been evaluated as candidate vaccines and may be included in next-generation vaccines. In areas where anaplasmosis is endemic and control measures are not applied, most calves become infected, resulting in endemic stability with
minimal clinical disease but high infection rates. See ANTIBIOTIC; VACCINATION.
Economic impact. Anaplasmosis continues to be a major disease of cattle in the tropical and subtropical world. Economic losses result from antibiotic control and vaccine costs, vector control, and production losses (such as decreased weight gain) from infected and recovered cattle and deaths. It is the most important tick-borne disease of cattle in the United States. See RICKETTSIOSES. Edmour F. Blouin Bibliography. A. F. Barbet, Recent developments in the molecular biology of anaplasmosis, Vet. Parasitol., 57:43–49, 1995; E. F. Blouin et al., Applications of a cell culture system for studying the interaction of Anaplasma marginale with tick cells, Animal Health Res. Rev., 3(2):57–68, 2002; J. de la Fuente et al., Molecular phylogeny and biogeography of North American isolates of Anaplasma marginale (Rickettsiaceae: Ehrlichiae), Vet. Parasitol., 97:65–76, 2001; D. M. French et al., Expression of Anaplasma marginale major surface protein 2 variants during persistent cyclic rickettsemia, Inf. Immun., 66(3):1200–1207, 1998; K. M. Kocan, Adaptations of the tick-borne pathogen, Anaplasma marginale, for survival in ticks and cattle, Exp. Appl. Acarol., 28:9–25, 2002; G. H. Palmer and T. F. McElwain, Molecular basis for vaccine development against anaplasmosis and babesiosis, Vet. Parasitol., 57:233–253, 1995.
Anapsida Formerly, a group of reptiles that was recognized in rank-based classifications. Amniotes were once divided into subclasses, based on the patterns of temporal fenestrae (openings in the skull roof posterior to the openings for the eyes). Subclass Synapsida included taxa that exhibit a single pair of temporal fenestrae (synapsid condition). Subclass Diapsida included taxa that feature two pairs of temporal fenestrae (diapsid condition) and those that are descended from ancestors that had them. The absence of temporal fenestrae is known as the anapsid condition. Mesosaurids, pareiasaurids, and captorhinids are examples of fossil reptiles that lack temporal fenestrae, whereas turtles are the only living reptiles that lack them; these reptiles were placed together in the subclass Anapsida. Certain Paleozoic reptiles (for example, millerettids) that exhibited synapsid-like temporal fenestrae but clearly did not belong in Synapsida (or Diapsida) were placed in Anapsida. See AMNIOTA; CAPTORHINIDA; CHELONIA; DIAPSIDA; MESOSAURIA; REPTILIA; SYNAPSIDA. Studies of early amniotes that use cladistic methodology, in which derived features are used to group taxa, revealed that the absence of temporal fenestrae is a primitive feature. Moreover, amniote taxa that had been placed in the subclass Anapsida in rank-based classifications were found to form a paraphyletic (artificial) group at the base of Reptilia in cladograms. Phylogenetic systematists refuse to recognize artificial groups with formal taxon names;
Anaspidacea but with the advent of phylogenetic nomenclature, attempts were made to conserve traditional group names by linking them with extant representatives, in order to adapt them as “clade names.” Thus, Anapsida was linked to a clade that included turtles. In the earliest cladistic studies of amniotes, turtles and captorhinids formed a clade, which was called Anapsida. Subsequent studies placed turtles in a clade that contained pareiasaurs, mesosaurs, and their close relatives, and the name Anapsida was transferred to that clade. That act generated confusion because captorhinids, which are regarded by paleontologists as exemplary fossil reptiles with anapsid skulls, were no longer to be regarded as anapsid reptiles. See ANIMAL SYSTEMATICS; PHYLOGENY; TAXONOMY. Recent discoveries have led taxonomists to advocate the abandonment of Anapsida as a taxonomic entity. Paleontological research indicates that synapsid-like temporal fenestrae arose at least four times in the clade of pareiasaurs, mesosaurs, and their close relatives. Morphological and molecular studies now place turtles in Diapsida, a corollary of which is that the anapsid condition of turtles can be interpreted as a derived feature of turtles. If turtles are truly nested within Diapsida, the name Anapsida can no longer be used as a clade name (because the phylogenetic definition for Anapsida used Diapsida as a reference taxon). Whereas it is no longer appropriate to use the term anapsid in a taxonomic context, it can still be used in a morphological sense to describe the skull of an amniote that lacks temporal fenestrae. Thus, the anapsid condition was common among Paleozoic amniotes, where it represented the retention of a primitive feature in those taxa. Conversely, the absence of temporal fenestrae in turtles is a derived feature of these reptiles (if the diapsid origin of turtles is correct). Sean Modesto Bibliography. M. J. Benton, Vertebrate Palaeontology, 3d ed., Blackwell, Oxford, 2004; R. L. Carroll, Vertebrate Paleontology and Evolution, W. H. Freeman, New York, 1988; A. S. Romer, Vertebrate Paleontology, 3d ed., University of Chicago Press, 1966.
Reconstruction of Rhyncholepis parvula from the Silurian of Norway. (Modified from A. Ritchie, 1980; reproduction by permission of Royal Society of Edinburgh, Trans. Roy. Soc. Edinburgh: Earth Sci., 92:263–323, 2002)
jawless fishes (such as Osteostraci, Heterostraci), and the trunk bears rod-shaped scales arranged in a chevronlike pattern in five variously inclined overlapping rows. The relatively fragile nature of this dermal skeleton means that well-preserved articulated fossils are rare but isolated scales and plates are fairly common. The anaspids are also characterized by a median dorsal row of peculiar, often hookshaped spines/scales. The tail is strongly hypocercal (the muscular lobe containing the notochord turns downward). Some articulated specimens have traces of paired ventrolateral fins, which vary in length from genus to genus (see illustration). Due to their close external similarities to living lampreys, anaspids have traditionally been regarded as close relatives or ancestors of the former. However, current views on the relationships of early vertebrates suggest that anaspids, together with other fossil jawless vertebrates, are more closely related to the jawed vertebrates (gnathostomes) than to lampreys. Their exact position, though, on the family tree is still a matter of considerable debate, mainly due to a lack of information regarding their internal anatomy. See ACANTHODII; HETEROSTRACI; JAWLESS VERTEBRATES; OSTEOSTRACI. Henning Blom Bibliography. H. Blom, T. M¨arss, and G. C. Miller, Silurian and earliest Devonian Birkeniid anaspids from the Northern Hemisphere, Trans. Roy. Soc. Edinburgh: Earth Sci., 92:263–323, 2002; P. Janvier, Early Vertebrates, Oxford Monogr. Geol. Geophys. 33, Oxford University Press, 1996.
Anaspidacea Anaspida A group of fossil jawless fishes, roughly similar in appearance to extant lampreys, known from Early Silurian–Early Devonian deposits (about 410– 430 million years ago) that represent marine, coastal environments. Anaspid fossils are known predominantly from the Northern Hemisphere. Most anaspids, such as Rhyncholepis, are small, rarely exceeding 15 cm (6 in.), with a rather slender and laterally compressed body. A slanting row of 7– 15 branchial (gill) openings is present between the head region and the trunk. Triradiate postbranchial spines at the ventral end of the branchial row characterize the group. The head of most anaspids is covered by a series of smaller plates, in contrast to the massive headshield seen in most other contemporary
An order of the crustacean superorder Syncarida, the most primitive of the Eumalacostraca. The families Anaspididae and Koonungidae have living representatives, while the Clarkecarididae are true fossils. The families Gampsonychidae and Palaeocarididae, formerly included in this order, now are assigned to a new order, Palaeocaridacea. The Anaspidacea arose during the Paleozoic Era, and the recent species are now restricted as living fossils to freshwater habitats in Tasmania and South Australia (not New Zealand). They exhibit a great adaptability to special habitats. The first thoracic somite is incorporated with cephalic tagmata. Anaspididae. Anaspides tasmaniae has a body which is similar to that of the amphipods. It is about 2 in. (5 cm) long and is found in mountain lakes and cold rivers at elevations of 2000–4000 ft
651
652
Anatomy, regional (600–1200 m) on Mount Wellington. Paranaspides lacustris, from the same region, is 1.2 in. (3 cm) long and resembles Mysis. Micrasides calmani is the smallest species and is found associated with Sphagnum. Anaspidites antiquus is from the Triassic of Australia. See AMPHIPODA. Koonungidae. Koonunga cursor is 0.8 in. (2 cm) long and is found in the mud bottom of temporary springs near Melbourne. The eyes are sessile and the first thoracic limb is modified for digging. See SYNCARIDA. Hans Jakobi Bibliography. R. D. Barnes, Invertebrate Zoology, 1968; H. K. Brooks, On the fossil Anaspidacea, with a revision of the classification of the Syncarida, Crustaceana, 4:239–242, 1962; S. P. Parker (ed.), Synopsis and Classification of Living Organisms, 2 vols., 1982; R. Siewing, Uber die Verwandtschaftsbeziehungen der Anaspidaceen, Zool. Anz. Suppl., 18:242-252, 1955; G. Smith, On the Anaspidacea, living and fossil, Quart. J. Microsc. Sci., 53:489-578, 1909.
elbow, ventral (anterior) medial (internal) posterior region
hand (palmar)
ventral brachial
parietal temporal orbital nasal oral neck, ventral (anterior) lateral (external) deltoid
medial brachial
sternal
axillary
mammary hypochrondriac epigastric umbilical mesogastric inguinal hypogastric pubic femoral, medial (internal) ventral (anterior) lateral (external)
Anatomy, regional The detailed study of the anatomy of a part or region of the body of an animal, most commonly applied to regional human anatomy. This is in contrast, but supplementary, to the study of organ systems, such as the cardiovascular, where all the structures pertaining to the system are studied in their continuity. To obtain the maximum information about the anatomy of an animal, both methods are used, as well as histological, embryological, and functional studies. The various regions of the body have particular interest for certain medical and other scientific specialists. For example, the obstetrician has a detailed knowedge of the perineal region, and the ophthalmologist of the eye. Regional anatomy is of great importance to the surgeon, and certain texts concerning it are labeled surgical anatomy. Closely related to regional anatomy is topographic anatomy, in which bony and soft tissue landmarks on the surface of the body are used to indicate the known location of deeper structures. An example is the point halfway between the umbilicus and the ventral prominence of the right hipbone; this marks the approximate location of the vermiform appendix and is known as McBurney’s point. There are many methods of dividing the body into regions for study, and one such means of classification is shown in the illustration. This system includes only externally visible areas; other systems would include special internal regions as well. The head, trunk, and extremities are the principal regions. Subdivision of these can be carried out, as illustrated, so that many distinct areas or especially vital regions are indicated. There is no end to the dividing and subdividing a specialist may do to make his task eventually less difficult. Regional anatomy is the natural and historical approach to the study of human and animal forms. The ancient scientists did not know much about the sys-
frontal
calcaneal
elbow lateral (external) forearm, palmar radial hand (dorsal) ungual patellar crural, lateral (external) ventral (anterior) medial (internal)
foot (dorsal)
Some of the major regions of the human body seen in ventral view. (After B. Anson, ed., Morris’ Human Anatomy, 12th ed., McGraw-Hill, 1966)
tems of the body except through conjecture. They approached the dissection of animals and cadavers in a manner intended to display the structures of some part, such as limb. The integration of such information led to the development of systematic anatomy. Thomas S. Parsons
Andalusite A nesosilicate mineral, composition Al2SiO5, crystallizing in the orthorhombic system. It occurs commonly in large, nearly square prismatic crystals. The variety chiastolite has inclusions of dark-colored carbonaceous material arranged in a regular manner. When these crystals are cut at right angles to the c axis, the inclusions form a cruciform pattern (Fig. 1). There is poor prismatic cleavage; the luster is vitreous and the color red, reddish-brown, olive-green, or bluish. Transparent crystals may show strong dichroism, appearing red in one direction and green in another in transmitted light. The specific gravity is 3.1– 3.2; hardness is 7.5 on Mohs scale, but may be less on the surface because of alteration. Occurrence and use. Andalusite was first described in Andalusia, Spain, and was named after this locality. It is found abundantly in the White Mountains near
Andesine
27 mm Fig. 1. Andalusite, variety chiastolite. Prismatic crystal specimens from Worchester County, Massachusetts. (American Museum of Natural History specimens)
Laws, California, where for many years it was mined for manufacture of spark plugs and other highly refractive porcelain. Chiastolite, in crystals largely altered to mica, is found in Lancaster and Sterling, Massachusetts. Water-worn pebbles of gem quality are found at Minas Gerais, Brazil. See SILICATE MINERALS. Cornelius S. Hurlbut, Jr. Aluminum silicate phase relations. The three polymorphs of Al2SiO5 are andalusite (Fig. 1), kyanite, and sillimanite. These minerals occur in medium- and high-grade metamorphic rocks, and andalusite and sillimanite occur rarely in granites. Laboratory determination of the pressure-temperature stability fields of these minerals has been difficult and remains a subject of controversy, but the phase diagram shown in Fig. 2 is generally accepted. It is consistent with thermochemical data and natural occurrence of temperature, °F 1000 1200 800
600
21
6 H2O pressure, kb
1400
kyanite sillimanite
4
14
7
2 andalusite
200
300
400
500
600
700
800
temperature, °C Fig. 2. Phase diagram for the polymorphs of aluminum silicate (Al2SiO5). 1 km = 0.6 mi; 1 kb = 102 MPa. (After M. J. Holdaway, Stability of andalusite and the aluminum silicate phase diagram, Amer. J. Sci., 271:97–131, 1971)
depth, km
400
the minerals. The univariant boundaries of the three stability fields meet at an invariant point located at 3.8 kilobars (380 megapascals) and 930◦F (500◦C). Grade (degree) of metamorphism is commonly a function of increasing temperature at some nearly constant or slowly changing pressure. Temperature variation is a reflection of changing temperature gradients (or heat sources) in the Earth, and pressure is an indication of thickness of overlying rock at the time of metamorphism. In pelitic (former shale) schists the grade or degree of metamorphism may be expressed by index minerals, such as the aluminum silicates. In regional metamorphic terranes the sequence of kyanite followed by sillimanite with increasing grade or temperature corresponds to pressures above 3.8 kbar (380 MPa) or depths greater than 8 mi (13 km) in the Earth’s crust. In other regional metamorphic terranes the sequence andalusite to sillimanite indicates shallower levels in the crust (Fig. 2). For the New England region of the United States, the sequence is kyanite to sillimanite in Vermont, western New Hampshire,western Massachusetts, and Connecticut, while in Maine, eastern New Hampshire, eastern Massachusetts, and Rhode Island the sequence is andalusite to sillimanite, all of which formed during Early Devonian metamorphism. This indicates that central and western New England were metamorphosed at greater pressure (deeper in the crust) than northern and eastern New England. Thus, natural occurrence of aluminum silicates combined with knowledge of their phase relationships is a useful aid in deducing extent of vertical movement and erosion occurring in the Earth’s crust after metamorphism. See KYANITE; METAMORPHISM; SILLIMANITE. M. J. Holdaway Bibliography. W. A. Deer, R. A. Howie, and J. Zussman, Rock-Forming Minerals, vol. 1a: Orthosilicates 2d ed., 1997; M. J. Holdaway, Stability of andalusite and the aluminum silicate phase diagram, Amer. J. Sci., 271:97–131, 1971; D. M. Kerrick, The Al2SiO5 Polymorphs, Reviews in Mineralogy, vol. 22, 1990; C. Klein, Manual of Mineralogy, 21st ed., 1993; E-an Zen, W. S. White, and J. B. Hadley (eds.), Studies of Appalachian Geology: Northern and Maritime, 1968.
Andesine A plagioclase feldspar with composition Ab70An30 to Ab50An50 (Ab = NaAlSi3O8; An = CaAl2Si2O8). Andesine occurs primarily in igneous rocks, often in a glassy matrix as small, chemically zoned, lathlike crystals known as microlites. The rock types may be called andesinites (if dominantly feldspar), andesites (see illus.), andesitic basalts (or olivine-bearing andesites, as in Hawaiian lava flows), or pyroxene-, hornblende- or biotite-andesites (depending on the second most dominant mineral—all are volcanic). See ANDESITE. The symmetry of andesine is triclinic, hardness on the Mohs scale 6, specific gravity 2.69, melting point ∼1210◦C (2210◦F). If quenched at very high
653
654
Andesite
2 cm Andesine grains with biotite in a specimen from Zoutpansberg, Transvaal, South Africa. Andesine is rarely found except as grains in igneous rocks.
temperatures, andesine has an albitelike structure, with aluminum (Al) and Silicon (Si) essentially disordered in the structural framework of the crystals. But, in the course of cooling, most natural andesines develop an Al-Si ordered structure with unusual, straininduced, periodic features detectable only by singlecrystal x-ray diffraction or by transmission electron microscopy. These are called e-plagioclases, based on the presence of non-Bragg diffraction maxima (ereflections) in x-ray or electron patterns. See ALBITE; CRYSTAL STRUCTURE. Calcic andesines and labradorites (Ab55An45– Ab40An60) may exsolve into two distinctly intergrown lamellar phases whose regularity of stacking produces beautiful interference colors like those in the feathers of a peacock. Polished specimens of this material are called spectrolite in the gem trade, and at some localities (notably eastern Finland) crystals up to 10 in. (25 cm) are mined by hand. Smaller crystals are made into cabochons for jewelry. They may be abundant enough in the host rock to be valued as a decorative stone. Hematite (Fe2O3) is present in rare oligoclases, and some andesines as micrometer-scale thin flakes, oriented parallel to certain structurally defined planes, producing brilliant reddish-gold reflections. These semiprecious stones are called aventurine or sunstone. See FELDSPAR; GEM; HEMATITE; IGNEOUS ROCKS; LABRADORITE. Paul H. Ribbe
South America, and the Taupo Volcanic Zone of New Zealand (Fig. 1) are andesitic. See LITHOSPHERE; PLATE TECTONICS; VOLCANO. Andesite volcanoes have been responsible for some of the most destructive, globally significant eruptions of historic and modern times. Examples include the 1815 eruption of Tambora and the 1883 eruption of Krakatau in Indonesia, the 1902 eruption of Mount Pelee in the West Indies, the 1980 eruption of Mount St. Helens in the United States, the 1991 eruption of Pinatubo in the Philippines, and the 1995–1998 eruption of the Soufriere Hills in the West Indies. The activity of most active andesite volcanoes is monitored continuously, but loss of life can still occur. For example, mud flows (lahars) from Nevado del Ruiz volcano in Colombia caused around 25,000 deaths in November 1985. Andesites are mostly dark-colored vesicular volcanic rocks which are typically porphyritic (containing larger crystals set in a fine groundmass). Phenocrysts (the larger crystals) comprise plagioclase; calcium-rich, calcium-poor pyroxene; and iron-titanium oxides set in a fine-grained, frequently glassy, groundmass. Some andesites contain phenocrysts of olivine, and some contain amphibole and biotite; these latter rocks generally contain more potassium. The porphyritic nature of andesites is derived from a complicated history of magmatic crystallization and evolution as the melts rise toward the surface from deep in the Earth. Phenocryst minerals commonly are strongly zoned and show evidence for disequilibrium during growth, consistent with an origin involving crystal fractionation and mixing processes. Andesites are readily classified in terms of their silicon dioxide (SiO2) content, between 53 and 63 wt %, and potassium oxide (K2O) content at a given SiO2 content (Fig. 2a; table). They can also be readily discriminated on a total alkali versus SiO2 diagram (the TAS diagram; Fig. 2b). Most andesite volcanoes erupt lavas and tephras (volcanic ash) which range in composition from basaltic andesite to dacite. Eruptions are often explosive, reflecting the relatively high water and gas content of the magmas. During an eruption, a column may
Andesite A typical volcanic rock erupted from a volcano associated with convergent plate boundaries. The process of subduction, which defines convergent plate boundaries, pushes oceanic lithosphere beneath either oceanic lithosphere or continental lithosphere. Andesites are the principal rocks forming the volcanoes of the “ring of fire,” the arcuate chains of volcanoes which rim the Pacific Ocean basin. The Marianas and Izu-Bonin islands, the islands of Japan, the Aleutian Islands, the Cascades Range of the northwest United States, the Andes mountain chain of
Fig. 1. Mount Ruapehu (2797 m; 9140 ft) in the North Island, New Zealand, erupting July 1996. This active andesite volcano is at the southern end of the Taupo Volcanic Zone continental margin arc system. Tephra (volcanic ash) is shown falling from the eruptive column to blanket the landscape to the north (left) of the volcano.
Andesite Chemical analyses (expressed in percent) of typical andesites from a continental margin arc (Mount Ruapehu, New Zealand) and an oceanic island arc (Raoul Island, Kermadec Islands, Southwest Pacific)
Raoul Percent
SiO2 TiO2 Al2 O3 Fe2 O3 MnO MgO CaO Na2 O K2 O P2 O5 LOI ∗ Total
59.44 0.67 16.92 6.37 0.10 4.03 6.79 3.38 1.70 0.12 0.27
57.27 1.23 14.10 13.31 0.23 2.93 7.52 2.69 0.62 0.19 ⫺0.47
59.1 0.70 15.8 6.6† 0.11 4.4 6.4 3.2 1.9 0.2
99.79
99.62
98.41
high K
Trace elements, ppm Sc V Cr Ni Cu Zn Ga Rb Sr Y Zr Nb Cs Ba La Ce Pr Nd Sm Eu Gd Tb Dy Ho Er Tm Yb Lu Hf Ta Pb Th U
16.7 152 51 25 49 68 17 59 274 19 130 4.54 5.13 339 13.54 29.24 3.8 14.92 3.49 1.04 3.87 0.57 3.29 0.69 1.87 0.30 2.03 0.31 4.26 0.81 14.53 5.47 1.63
∗ Loss on ignition ( ≈ volatile † All Fe as Fe O . 2 3
41 294 4 4 102 122 13 8.5 188 37 75 0.82 0.35 202 4.4 13.0
3.15
22 131 119 51 24 73 16 58 325 20 123 12 2.6 390 18 42 5.0 20 3.9 1.2 3.6 0.56 3.5 0.76 2.2
3.95 0.62 2.37 0.066 2.7 0.563 0.233
2.0 0.33 3.7 1.1 12.6 5.6 1.4
11.7 3.84 1.35 4.61 0.93 5.41
rhyolite
4.0
content).
rise tens of kilometers above the volcano, leading to stratospheric dispersion of tephra and volcanic gases and generation of pyroclastic flows (for example, Pinatubo). Pyroclastic flows are a particular feature of andesite-type volcanism and are among the most dangerous of volcanic hazards. Indeed, it was pyroclastic flows from Vesuvius which overwhelmed Pompeii in A.D. 79; and at Mount Pelee on the island of Martinique in 1902 the pyroclastic flow engulfed the town of St. Pierre, killing around 28,000 people. See BASALT; LAVA; PYROCLASTIC ROCKS. Less commonly, rocks of andesite and basaltic-
3.5 weight percent oxide (K2O)
Ruapehu
andesite composition are associated with sites of intraplate volcanism, such as Iceland, Gal´apagos, and Hawaii, and are unrelated to subduction or orogenic processes. These rocks have basaltic-andesite to andesite composition (that is, in the composition range of 53–63% SiO2) and have been called icelandite and hawaiite. They have formed in response to fractional crystallization of basaltic magma. It is generally recognized that most andesites cannot be direct partial melts of peridotitic mantle. As such, andesites are evolved magmas and the end product of a plethora of processes, including crystal
dacite andesite
3.0 2.5 2.0
basaltic andesite
medium K
basalt
1.5 1.0
low K
0.5 52
56
60 64 weight percent oxide (SiO2)
68
72
(a) 8 7 weight percent oxide (Na2O + K2O)
Component
Average continental crust
655
6 5 4 3 2 basaltic andesite
1 0 50
52
54
andesite
56 58 weight percent oxide (SiO2)
60
62
64
(b) Fig. 2. Igneous rock and andesite classification. (a) Igneous rocks, in terms of SiO2 versus K2O. Note that a continuum exists between basalt, basaltic andesite, andesite, dacite, and rhyolite, and that it is possible to subdivide igneous rocks into low-, medium-, and high-K varieties. The data shown are andesites from Mount Ruapehu, which classify in this diagram as medium-K basaltic andesites and andesites. (b) Andesitic rocks, in terms of SiO2 versus (Na2O + K2O), showing a part of the total alkali versus silica (TAS) diagram. (After M. J. Le Bas et al., A chemical classification of volcanic rocks based on the total alkali-silica diagram, J. Petrol., 27:745–750, 1986)
656
Andesite ARC VOLCANO storage, fractionation, and mixing
68 mi (110 km)
ARC LITHOSPHERE
storage, fractionation, and mixing subduction angle
w (H ate 2O r )
melt aggregation and fractional crystallization MANTLE WEDGE
AB f SL s o IC ed ER ndr H P hu r OS to yea H T s r LI ten pe : G N e rs TI rat ete C n m U io lli BD ct mi SU bdu su zone of partial melting
Fig. 3. Schematic cross section through a typical subduction system. The major components of a convergent plate boundary are the arc lithosphere, the subducting lithospheric slab, and the mantle wedge. Note that arc volcanoes overlie the subducting slab by around 110 km (68 mi), that the volume of the mantle wedge is determined by arc lithosphere thickness and slab dip, and that the mantle wedge is the principal site of magma production. After magmas form by partial melting, they can evolve further by fractional crystallization, mixing, and storage in reservoirs over a range of depths.
fractionation, magma mixing, mingling, assimilation, and storage, which have acted on primary basaltic magmas produced by partial melting in the mantle wedge. See MAGMA. In more detail, the process of subduction recycles oceanic lithosphere back into the deep Earth; and as such, it complements the process of sea-floor spreading which generates new basalt crust and oceanic lithosphere at the mid-ocean ridges. The zone along which subduction takes place is demarcated at the Earth’s surface by a deep oceanic trench and within the Earth by an increase of earthquake intensity along an inclined plane, the Wadati-Benioff Zone. This zone derives from the fact that cold, brittle material (the subducting oceanic lithosphere) is pushing into warmer upper mantle, and the earthquakes reflect the contrasting physical properties between the two mediums. The subduction process takes place at rates varying from a few tens to a few hundreds of millimeters per year and at angles of around 20◦ (shallow subduction) to nearly 90◦ (steep subduction). Typically, andesite volcanoes are located about 110 km (68 mi) above the surface of the subducting plate
(Fig. 3). Between the surface of the subducting plate, which geologists refer to as the subducting slab, and the overlying plate is the mantle wedge, where basaltic melts form by the process of partial melting. Partial melting involves the generation of melt by dissolution of minerals, leaving a residual, more refractory, mineral assemblage. The melts aggregate and then move toward the surface by virtue of their lower density and viscosity. The melting occurs because the subducting slab, which may have spent many millions of years on the ocean floor as oceanic crust, contains minerals with seawater locked into their structure. When the minerals are subjected to higher pressure and temperature as subduction proceeds, chemical reactions in the minerals in the slab lead to loss of water and other volatile constituents which diffuse upward into the mantle wedge, together with an array of readily soluble elements such as potassium, rubidium, cesium, barium, and uranium. The mantle wedge becomes relatively enriched in these elements. The introduction of water is especially significant as it depresses the melting temperature and facilitates partial melting among constituent mineral grains. Importantly, the mantle wedge is recognized as the source of most arc-related magmas. On rare occasions, where subducting oceanic lithosphere is young, is still hot, and is subducting rapidly, the slab itself may experience partial melting, generating a unique suite of high-silica, low-potassium rocks of andesite composition, called adakites (after Adak Island in the Aleutian island arc chain). Still another unique suite of andesitic rocks, called boninites, which have high-magnesium and typically low-titanium contents, lacking plagioclase phenocrysts but with magnesium-rich, calcium-poor pyroxene phenocrysts, are associated with some suites of island arc rocks (for example, the Bonin island arc, which is the type area), where watersaturated melting of depleted mantle peridotite produced partial melts of high-silica, highmagnesium content. See SUBDUCTION ZONES. The bulk composition of typical andesites is similar to that of continental crust (see table), contrasting with oceanic crust which is basaltic. Over geological time, subduction and the resulting andesite volcanism has proved a means of generating new continental crust, possibly at continental margin arc systems. John Gamble Bibliography. P. W. Francis, Volcanoes: A Planetary Perspective, Oxford University Press, 1993; J. A. Gamble et al., A fifty year perspective of magmatic evolution on Ruapehu Volcano, New Zealand: Verification of open system behaviour in an arc volcano, Earth Planet. Sci. Lett., 170:301– 314, 1999; J. B. Gill, Orogenic Andesites and Plate Tectonics, Springer-Verlag, Berlin, 1981; M. J. Le Bas et al., A chemical classification of volcanic rocks based on the total alkali-silica diagram, J. Petrol., 27:745–750, 1986; A. R. McBirney, Igneous Petrology, Jones and Bartlett, Boston, 1993; R. L. Rudnick, Making continental crust, Nature, 378:571–578, 1995; R. S. Thorpe (ed.), Andesites: Orogenic Andesites and Related Rocks, Wiley, Chichester, 1982;
Androgens M. Wilson, Igneous Petrogenesis: A Global Tectonic Approach, Unwin Hyman, London, 1989.
Andreaeopsida A class of the plant division Bryophyta, containing plants commonly called granite mosses. The class consists of one order and two families, the Andreaeobryaceae and the Andreaeaceae. The large tapered foot and short, massive seta of the Andreaeobryaceae show possible linkage to the Bryopsida, but the differences seem more fundamental and more significant than the similarities. The small, brittle plants grow as perennials on rocks in montane regions in dark red, red-brown, or blackish tufts. The stems are erect, simple or forked, and grow from a single apical cell with three cutting faces; they consist of a homogeneous tissue of thick-walled cells. The rhizoids are multicellular and filamentous or platelike. The leaves are spirally arranged, of various shapes and with or without a midrib. The leaf cells are thick-walled and often papillose. Paraphyses are present in both male and female inflorescences. The sex organs are superficial in origin. The antheridia are long-ellipsoid and stalked. The archegonia are flask-shaped. The terminal sporophytes consist of a small foot and an ellipsoid capsule elevated on an extension of the gametophytic axis (see illus.), the pseudopodium (or, in Andreaeobryum, they consist of a well-developed foot, a short seta and no pseudopodium). The capsule dehisces by 4– 8 vertical slits. The slender columella is surrounded
sporophytes
Granite moss (Andreaea ruperstris), showing gametophyte with sporophytes. (After G. M. Smith, Cryptogamic Botany, vol. 2, 2d ed., McGraw-Hill, 1955)
(and overarched) by sporogenous tissue of endothecial origin. The wall consists of solid tissue. Stomata, peristomes, and elaters are lacking. The spores begin division within the spore wall before dehiscence of the capsule. The protonema is a branched filament producing numerous buds; erect filaments may produce at their tips leaflike plates. The calyptra is small and mitrate or, in Andreaeobryum, larger and campanulate-mitrate. See ANTHOCEROTOPSIDA; BRYOPHYTA; BRYOPSIDA; HEPATICOPSIDA; SPHAGNOPSIDA. Howard Crum Bibliography. S. P. Parker (ed.), Synopsis and Classification of Living Organisms, 2 vols., 1982; W. Schultze-Motel, Monographie der Laubmoosgattung Andreaea, I: Die costaten Arten, Willdenowia, 6:25– 110, 1970; W. C. Steere and B. M. Murray, Andreaeobryum macrosporum, a new genus and species of Musci from northern Alaska and Canada, Phytologia, 33:407–410, 1976.
Androgens A class of steroid hormones produced in the testes and adrenal cortex which act to regulate masculine secondary sexual characteristics. The major naturally occurring androgens, such as testosterone, are C-19 steroids (that is, they all contain 19 carbon atoms) with a hydroxyl or ketone group at the C-3 and C-17 positions (see testosterone in Fig. 1). Reproductive functions. Androgens play several important roles related to male reproductive functions. These steroids are responsible for the development of the entire reproductive tract, including the epididymis, seminal vesicles, and vas deferens, in the developing male embryo as well as the external genitalia, including the penis. At the onset of puberty, androgen secretion increases dramatically; once testosterone levels reach a plateau, they remain fairly constant until a man reaches the age of 70–80, after which they slowly decline. During puberty, these androgenic hormones are essential for development of the masculine secondary sex characteristics, which include a deepening of the voice, growth of a beard and coarse body hair, growth of the penis, and increased development of muscles. They are also required for development of the male sexual drive, or libido; initiation and maintenance of spermatogenesis, or sperm production; and the growth and function of accessory sex glands, including the prostate gland. See PROSTATE GLAND; SPERMATOGENESIS. Although the interstitial cells (Leydig cells) in the male testes secrete testosterone, other less potent androgens, including dehydroepiandrosterone (DHEA) and androstenedione, are secreted by the cortex of the adrenal glands in both males and females. In females these adrenal androgens modulate the physiological function of reproductive and sexual organs including the ovaries, vagina, clitoris, and mammary glands. They are also responsible for growth of hair in the armpits and appear to be required for sexual arousal of the female. Overproduction of these adrenal androgens in a female can cause
657
658
Androgens Acetate Cholesterol 21 CH3 1 20 C
CH3 O
C
O OH
O
2
∆5 pathway OH
1 4
6 OH
Pregnenolone 3
CH3 C
∆4 pathway
17α-Hydroxypregnenolone
OH
Dehydroepiandrosterone
CH3
3
C
O 2
2
OH
O OH
A O
OH
O 17α-Hydroxyprogesterone
Progesterone
Androstenediol 3
O
OH
18 19 O
O
17
3
OH 4
HO Testosterone
Androstenedione
Estradiol
5
Enzymes: 1 20,22-desmolase 2
17α-hydroxylase
3
3β-hydroxysteroid dehydrogenase
4
aromatase 5α-reductase
5
OH
O
H Dihydrotestosterone
Fig. 1. Major pathways of testicular androgen and estrogen synthesis.
masculinization, which includes the growth of coarse body hair. Although the synthesis and secretion of adrenal androgens in both males and females is stimulated by adrenocorticotropic hormone, testosterone secretion by the Leydig cells of the male testes is specifically stimulated by luteinizing hormone. Both adrenocorticotropic hormone and luteinizing hormone are tropic hormones secreted by the anterior pituitary gland. See ADENOHYPOPHYSIS HORMONE; ADRENAL CORTEX; OVARY; REPRODUCTIVE SYSTEM; TESTIS. Nonreproductive functions. Androgens also produce nonreproductive effects that include increased secretory activity of the sebaceous glands in a boy’s skin at the time of puberty and fusion of the epiphyseal growth plate in the long bones, preventing an adult male from continuing to grow taller. DHEA and its sulfated ester, DHEA(S), have been reported to exert a number of potential health benefits that range from antiaging and anticancer effects to enhanced weight loss and stimulation of the immune system. However, the use of synthetic androgens, often referred to as anabolic steroids, to enhance muscle mass and athletic ability can have deleterious
consequences. In a teenage boy, these undesirable side effects include not only acne and breast development but also a shutting down of bone growth. In adult males, overuse of these anabolic steroids can result in: infertility (decreased sperm production); liver damage; acceleration of the growth of prostate tumors; and the development of liver and brain tumors. Biosynthesis. Like all steroid hormones, testosterone is synthesized from cholesterol, which in turn is either synthesized from acetate (Fig. 1) or derived from serum lipoproteins. The rate-limiting step for testosterone synthesis is the side-chain cleavage of cholesterol to form pregnenolone. The desmolase enzyme, located within the mitochondria of Leydig cells, catalyzes this key reaction. Once pregnenolone is generated, it is converted to testosterone via two different pathways. The 5 pathway is the major pathway in humans; in this pathway side-chain cleavage (carbons 20 and 21) and reduction of the 17-keto group occur before oxidation of the first ring (A ring) of the steroid nucleus. In the 4 pathway this sequence is required. The enzymes that catalyze these steroid modifications are identical for both pathways
Androgens Conversion to estradiol. Testosterone can also be converted to estradiol via an enzyme that catalyzes the aromatization (introduction of double bonds) of the A ring of the steroid (Fig. 1). This reaction occurs primarily within the testicular Sertoli cells when they are stimulated by another anterior pituitary tropic hormone, follicle stimulating hormone (FSH). Similar amounts of estradiol can also be generated from testosterone within adipose tissue and the central nervous system. Although estradiol is a wellcharacterized female reproductive hormone, normal concentrations are also required for spermatogenesis, as well as regulation of bone density in males. Overproduction of estradiol in obese men can lead to feminization, including enlargement of the breasts. See ESTROGEN. Catabolism. The metabolic breakdown, or catabolism, of testosterone takes place fairly rapidly in the liver. The major metabolites that are produced are 17-ketosteroids, generated from either androstenedione or dihydrotestosterone (Fig. 2). The major end product of this metabolic pathway is androsterone, which is excreted in the urine in the free form of conjugated with glucuronate or sulfate. DHEA, the predominant androgen secreted by the adrenal glands, is the major source of these excreted 17-ketosteroids. Other metabolites of testosterone that are found in small amounts in urine include
and are located in the smooth endoplasmic reticulum of testosterone-secreting cells. Once the synthesized testosterone has been released into the bloodstream, it circulates throughout the body bound to proteins that are synthesized in the liver. Most of the secreted testosterone binds with high affinity to a serum transport protein known as sex steroid-binding globulin. Roughly 40% of the circulating testosterone is bound with low affinity to serum albumin. Although only 1–3% of the circulating testosterone is unbound, it is this free form of testosterone that is biologically active. See CHOLESTEROL; STEROID. Further reduction of testosterone to generate the more potent androgen dihydrotestosterone (DHT) occurs in some androgen-responsive tissues, including cells in the prostate gland, that express the 5α-reductase enzyme (Fig. 1). This potent androgen binds to intracellular androgen receptors more tightly, or with higher affinity, than testosterone itself. Once either testosterone or dihydrotestosterone has bound to these receptor proteins located within target cells, these hormone-receptor complexes interact with other molecules (coactivators and corepressors) and bind directly to specific nuclear deoxyribonucleic acid (DNA) sequences. This nuclear binding subsequently regulates the expression or transcription of androgen-responsive genes. See GENE.
OH
O Testosterone 1
3
OH
OH
OH 5
2
O
HO Estradiol
HO
H
Dihydrotestosterone Metabolites (less potent)
2
4
3
O Dehydroepiandrosterone
O
O
O
O
HO
H Androstanediol
5
O Androstenedione
H Androstanedione
Enzymes: 1
aromatase
4
3β-hydroxysteroid dehydrogenase isomerase
2
17β-hydroxysteroid dehydrogenase
5
3β-ketoreductase
3
5α-reductase
Fig. 2. Catabolism of androgens to generate active metabolites or inactive 17-ketosteroids.
17 HO
H Androsterone
659
Andromeda androstanediol, which is formed by the reduction of the 3-keto group of dihydrotestosterone, and estrogen metabolites (estrone and estriol) which are formed after testosterone has been converted to estradiol. Thomas J. Schmidt Bibliography. C. J. Bagatelli and W. J. Bremner, Androgens in Health and Disease, Humana Press, Totowa, NJ, 2003; R. M. Berne et al., Physiology, 5th ed., Mosby, St. Louis, 2003; C. Chang, Androgens and Androgen Receptor: Mechanisms, Functions, and Clinical Applications, Kluwer Academic Publishers, Boston, 2002.
can be seen with the unaided eye, since the glow that comes from its core can be seen even though it is 2.4 million light-years away. Only photographs show the galaxy’s extent. See ANDROMEDA GALAXY. The modern boundaries of the 88 constellations, including this one, were defined by the International Astronomical Union in 1928. See CONSTELLATION. Jay M. Pasachoff
Andromeda Galaxy The giant spiral galaxy nearest to Earth. The Andromeda Galaxy is seen in the constellation of Andromeda, which is visible from Earth in the autumn skies. The galaxy is also known as M31 and NGC 224, based on its entries in two widely used catalogs. See ASTRONOMICAL CATALOG; CONSTELLATION; MESSIER CATALOG. The Andromeda Galaxy can be seen without a telescope, one of a very few external galaxies that are sufficiently close and bright. Visually through a telescope, it appears as a large, faint, nearly featureless object. But photography, especially with modern electronic detectors, shows that it is an immense system of stars and other material, arranged in a nearly edge-on disk and large, spherical, faint halo (Fig. 1). It has been mapped in great detail at many wavelengths, including radio, infrared, visual, ultraviolet, and x-ray.
Andromeda A prominently located constellation in the northern sky (see illustration), named for the daughter of Cassiopeia in Greek mythology: When Cassiopeia bragged that her daughter Andromeda was more beautiful than Poseidon’s daughters, the Nereids, Poseidon created Cetus, the sea monster. (In some versions of the myth, Cassiopeia boasted of her own beauty.) The situation required Andromeda’s sacrifice. However, Perseus saved Andromeda by showing Cetus the head of Medusa, turning Cetus to stone. See CASSIOPEIA; PERSEUS. The Great Galaxy in Andromeda (M31) is the nearest spiral galaxy to the Earth, and along with the galaxy M33 in Triangulum is the farthest object that
3h
right ascension 1h
2h
0h
23h Magnitudes:
+50°
PERSEUS 2.0–2.5
CASSIOPEIA
2.5–3.0 3.0–3.5
LACERTA 51
Almaak γ
+40°
ω
χ
τ
ψ
ξ
4.5–5.0 5.0–5.5 5.5–6.0
ι υ
ο
Variable
M31
Galaxy
ν ϑ
R ρ
Mirach
β
ANDROMEDA
σ
π
TRIANGULUM
δ ε
ARIES +20°
λ κ
µ
+30°
3.5–4.0 4.0–4.5
ϕ
declination
660
PISCES
Alpheratz
α
PEGASUS
ζ η
◦
Modern boundaries of the constellation Andromeda, daughter of Cassiopeia. The celestial equator is 0 of declination, which corresponds to celestial latitude. Right ascension corresponds to celestial longitude, with each hour of right ascension ◦ representing 15 of arc. Apparent brightness of stars is shown with dot sizes to illustrate the magnitude scale, where the brightest stars in the sky are 0th magnitude or brighter and the faintest stars that can be seen with the unaided eye at a dark site are 6th magnitude. (Wil Tirion)
Andromeda Galaxy
Fig. 1. Photographic image of the Andromeda Galaxy obtained with the 1.2-m Schmidt telescope of the Palomar Observatory. (Courtesy of Paul Hodge)
The Andromeda Galaxy is one of 36 galaxies that are clustered together in the Local Group. It and the Milky Way Galaxy are the largest members and are two of only three spiral galaxy members. Its distance can be measured in a number of ways, most accurately by the use of the Cepheid variable period-luminosity law. The distance is found to be 2,500,000 light-years (1.5 × 1019 mi or 2.4 × 1019 km). See CEPHEIDS; LOCAL GROUP. Size. The main body of the Andromeda Galaxy is about 3◦ across in the sky, corresponding to about 135,000 light-years (8 × 1017 mi or 1.3 × 1018 km). However, the hydrogen gas in the galaxy’s disk can be detected to almost twice that size, and the darkmatter halo, the extent of which is only roughly measured, is even larger. (Dark matter is material of an unknown nature that can be detected only by its gravitational influence on visible matter. It appears to be the major constituent of galaxies and galaxy clusters.) Contents. There are a few hundred billion stars (a few times 1011) in the Andromeda Galaxy, most of them apparently quite similar in their properties to the stars in the Milky Way Galaxy. The Andromeda Galaxy also is similar in containing hundreds of star clusters, both globular clusters (Fig. 2) and smaller, young clusters, many of which have been studied in detail by the Hubble Space Telescope. See HUBBLE SPACE TELESCOPE; STAR; STAR CLUSTERS. Between the stars and clusters there is an interstellar medium of gas, mostly atomic hydrogen, detected by radio telescopes. The gas is very tenuous except in certain areas where it is dense enough to condense into new stars. Star-forming regions are made up of both gas and dust, the gas in these places
661
being largely molecular hydrogen, as well as other molecules. The inner parts of the Andromeda Galaxy are relatively free of atomic hydrogen, but contain much cool molecular gas. Currently star formation is concentrated at intermediate distances from the center, where the spiral arms are most conspicuous. The total mass of gas is a few billion times the mass of the Sun. See INTERSTELLAR MATTER; RADIO ASTRONOMY. There is also dust located throughout the disk of the galaxy, visible both at optical and at infrared wavelengths. The northwest side of Andromeda is heavily obscured by this dust, as seen in telescopic images. The total mass of dust in the galaxy is about 1/100 the mass of gas. Early results from space-borne x-ray telescopes searching for x-ray sources in the Andromeda Galaxy turned up a few dozen sources, mostly located in the inner parts. Searches with the Chandra X-ray Observatory have pinpointed hundreds of x-ray sources, many of which correspond in position with globular clusters. Most of them are neutron stars like those found in Milky Way clusters. See CHANDRA X-RAY OBSERVATORY; NEUTRON STAR; X-RAY ASTRONOMY. Structure. The nucleus of the Andromeda Galaxy is unusual. As discovered by the Hubble Space Telescope, it is a double nucleus, the two components being too close together to have been resolved by ground-based telescopes. The two revolve around each other, and one contains a black hole of several million solar masses. See BLACK HOLE. Surrounding the nucleus is an ellipsoidal bulge, made up of old and intermediate-age stars.
Fig. 2. Globular star cluster G1 in the Andromeda Galaxy, imaged by the Hubble Space Telescope. (NASA)
662
Anechoic chamber Enveloping the whole disk and extending even farther out than the bulk of the disk stars is a nearly spherical halo, primarily containing thinly spaced old stars. However, unlike the halo of the Milky Way Galaxy, that of the Andromeda Galaxy is not mainly made up of ancient stars with low heavy-element abundances. The origin and evolution of the halo differs in some unexplained way from those of the Milky Way Galaxy. Paul Hodge Bibliography. P. Hodge, The Andromeda Galaxy, Kluwer Academic, 1992; P. Hodge, Atlas of the Andromeda Galaxy, University of Washington Press, 1981; T. R. Lauer et al., Planetary camera observations of the double nucleus of M31, Astron. J., 106: 1436–1447, 1993; B. F. Williams, Recent star formation history of the M31 disk, Astron. J., 126:1312– 1325, 2003.
Anechoic chamber A room whose boundaries absorb effectively all the waves incident on them, thereby providing free-field conditions. A free field is a field whose boundaries exert negligible effect on the incident waves. In practice, it is a field in which the effects of the boundaries are negligible over the frequency range of interest. Acoustic chambers and radio-frequency and microwave chambers will be discussed. Acoustic chambers. An acoustic anechoic chamber is a room in which essentially an acoustic free field exists. It is sometimes referred to as a free-field or dead room. The word anechoic is derived from the Greek, meaning “without echo.” However, a room that is free from echoes (in the usual sense of the word) is not necessarily an anechoic room: in addition, there must be no significant reflections of sound waves from the boundaries of the room. Free-field conditions can be approximated when the absorption by the boundaries of the room approaches 100%. To reduce sound reflected by the
Interior of an acoustic anechoic chamber. The room’s special construction eliminates 99.9% of reflected sound.
boundaries to a minimum, the absorption coefficient must be very high and the surface areas of the boundaries should be large. The sound absorptive material usually installed in such rooms consists of glass fibers held together with a suitable binder. In order to achieve large surface area, a wall construction (see illus.) is used that includes wedges of sound absorptive material, the base of which is usually between 8 × 8 in. (20 × 20 cm) and 8 × 24 in. (20 × 60 cm), and the length of which is usually 3 to 5 ft (0.9 to 1.5 m). These wedges resemble stalagmites and stalactites and absorb about 99% of incident sound energy over most of the audiofrequency ranges. A horizontal net of thin steel cables just above floor wedges permits walking in the room. See REFLECTION OF SOUND; SOUND ABSORPCyril M. Harris TION. Radio-frequency and microwave chambers. The radio-frequency or microwave anechoic chamber is a shielded (screened) enclosure in which internal reflections of electromagnetic waves are reduced to an absolute minimum, thus providing a zone that is free from radio noise, broadcast transmissions, and other extraneous signals, and that simulates freespace conditions. Its main use is to provide a controlled and well-defined environment for the testing of electronic equipment. Shielded or screened enclosures may vary in size from small metallic boxes that are intended for the protection of circuit elements to large rooms capable of housing vehicles or aircraft. The purpose of such enclosures is to isolate sources of electromagnetic energy from electronic, radio, or other sensitive equipment which might suffer degradation of performance as a result of interaction with transmitters, radio-frequency heating equipment, and so forth. The provision of a radio-frequency quiet zone is an essential feature of many applications in electromagnetic tests and investigations, for example, the determination of antenna characteristics and the measurement of low-level emissions during electromagnetic compatibility assessments. See ELECTRICAL SHIELDING; ELECTROMAGNETIC COMPATIBILITY; MICROWAVE FREE-FIELD STANDARDS. The mechanism by which shielded enclosures function is a combination of absorption as the electromagnetic wave progresses through the material of the shield and reflection of the wave at the interface. The principle of reciprocity applies, and reflection occurs at the internal as well as at the external interface. In most practical shielded enclosures any radio-frequency signals generated internally are reflected and the bulk of the energy remains within the enclosure. At certain frequencies where the dimensions of the enclosure are comparable with half a wavelength, reinforcement of reflected waves take place and the enclosure becomes a resonant cavity. The internal conditions are thus somewhat different from free space, and it is the attempt to reproduce free-space or open-field conditions within a shielded enclosure which has led to the development of resistive material to line the enclosure walls to absorb rather than reflect the
Anemia radio-frequency energy and hence create an anechoic chamber. See ABSORPTION OF ELECTROMAGNETIC RADIATION; CAVITY RESONATOR; RECIPROCITY PRINCIPLE; REFLECTION OF ELECTROMAGNETIC RADIATION. Theoretical investigations indicate that reflectionfree conditions can be established at discrete frequencies by provision of sheets of resistive material that have a surface impedance of 120π (377) ohms per square, that is, the free-space or plane-wave impedance, and are spaced at one-quarter wavelength from the reflecting surfaces. Extension of this principle to achieve broadband coverage and hence an anechoic enclosure capable of use over a wide frequency range led to the development of various types of material having the appropriate resistive, lossy characteristics. Such material must have good basic dielectric properties, with a relative permittivity near unity to minimize reflection at the interface, and to incorporate increasing resistive loss through the material. In practice, foam plastics are ideal for the purpose, and these can be readily loaded with carbon, ferrite, or other conducting particles on a graded or layer basis to provide the increasing loss. Furthermore, the use of cones or pyramids of the material assists considerably in the transition from minimal reflection at the apex to full absorption at the base where the resistive loading is at a maximum. The lining of the walls and ceiling of a shielded enclosure with such materials provides for absorption of any internal signals and therefore anechoic or minimal reflection conditions simulating propagation in free space. See DIELECTRIC MATERIALS; ELECTROMAGNETIC RADIATION; PERMITTIVITY. In the microwave band, that is, at frequencies greater than 1 GHz, the wavelength is less than 1 ft (30 cm), and hence a lining thickness of only 1 or 2 in. (2.5 or 5 cm) provides for a reflectivity better than −40 dB or less than 1%. Enclosures that are used for test and measurement purposes have walls and ceilings with conical or pyramidal linings up to 7 ft (2 m) thick and are capable of simulating freespace conditions with reflections less than 10% over a wide frequency range extending down to 50 MHz. Special precautions are taken with the installation of lossy material on the floor, because the latter generally needs to be mechanically strong in addition to having desirable radio-frequency characteristics. See RADAR-ABSORBING MATERIALS. G. A. Jackson Bibliography. P. A. Chatterton and M. A. Houlden, EMC: Electromagnetic Theory in Practical Design, 1992; C. M. Harris, Handbook of Acoustical Measurements and Noise Control, 3d ed., 1991.
Anemia A reduction in the total quantity of hemoglobin or of red blood cells (erythrocytes) in the circulation. Because it generally is impractical to measure the total quantity, measures of concentration are used instead. Hemoglobin is contained in red blood cells, which are suspended in plasma, the liquid compo-
nent of blood. Therefore, concentration is affected not only by quantities of hemoglobin and red blood cells but also by plasma volume. Thus, the apparent anemia found in many women in the third trimester of pregnancy is not really anemia at all: the red cell mass is actually increased, but the plasma volume is expanded even more. In other words, hemodilution is present. Conversely, in dehydration and other circumstances of hemoconcentration, the plasma volume is reduced, thereby tending to mask anemia. See BLOOD; HEMOGLOBIN. The three measures of concentration most often employed are the hemoglobin, the red blood cell count, and the volume of packed red cells. Hemoglobin is determined photocolorimetrically. The red blood cell count is determined by automated optical-electronic devices; until recently it was determined by microscopic observation of a counting chamber—a less accurate method. The volume of packed red cells, often called the hematocrit after the instrument originally used in its measurement, is a measure of the proportion of volume that red cells occupy in a specimen of whole blood, and is determined by centrifugal methods, or by computation from light scatter or displacement volume of individual cells during automated enumeration procedures. When total red cell mass must be known, it is measured by removing some of the patient’s red cells and tagging them with the radioisotope chromium-51, reinjecting them, and then computing the red cell mass from the resulting dilution of the tagged red cells by the unlabeled red cells in the individual’s bloodstream. For diagnostic and therapeutic purposes, anemia is categorized according to the average volume of the individual red cells (mean corpuscular volume) and the concentration of hemoglobin within the red cells (mean corpuscular hemoglobin concentration). General considerations. In a group of healthy individuals, the values for hemoglobin, red cell count, and hematocrit approximate a “normal” distribution. Values that are less than 2.5 standard deviations below the mean are indicative of anemia if other clinical factors do not indicate a condition of hemodilution. The mean values are greater for adult males than for adult females, and greater for adults than children. Lower atmospheric oxygen tension at higher altitudes results in higher mean values for hemoglobin, red cells, and hematocrit in healthy individuals living under these conditions; hence, anemia would be defined at a higher value. Red blood cells have a finite life-span, averaging 120 days in the circulation in healthy persons, but the span may be shortened in certain diseases. Every day, about 45 billion aged red cells are removed from each liter of blood, or a total of 2 × 1011 red cells per day in the adult of average size. In the state of health, the rate of production of new red blood cells (erythropoiesis) equals the rate of removal of senescent red blood cells. Anemia occurs if the rate of erythropoiesis is reduced below normal. It also occurs if hemorrhage or destruction (hemolysis) of red blood cells within the body increases the rate of loss of
663
664
Anemia erythrocytes and the rate of erythropoiesis does not increase enough to compensate. Mechanisms. Red blood cell production may be affected by several different mechanisms. Erythropoietin, a growth factor produced by healthy kidneys, stimulates the bone marrow to produce more erythroblasts and accelerates their maturation. An intact bone marrow can respond by increasing production of red blood cells by about sevenfold after a few days or weeks. The proliferative response of the bone marrow to anemia may be defective or absent if the production of erythropoietin is diminished or absent, as occurs in chronic renal disease and some endocrine disorders. Erythropoietin produced synthetically has become available on prescription to overcome that deficiency and thereby resolve that specific type of anemia. See KIDNEY. Ineffective erythropoiesis results when, in the intact, stimulated bone marrow, red blood cell precursors either fail to mature, or die in the bone marrow prior to their delivery to the circulation as erythrocytes. Aplastic anemia occurs when the bone marrow stem cells, that give rise to precursors of erythroblasts, are markedly diminished in number or may respond inadequately to erythropoietin. This condition occurs when the bone marrow is adversely affected by certain chemicals or autoantibodies; is injured by irradiation; atrophies, or is replaced by fat; is replaced by fibrous (scar) tissue; or is infiltrated by cancer cells. Drugs and other chemicals that have been associated with aplastic anemia include those whose effects are dose-related and are an extension of their intended purpose, for example, some of the drugs used in cancer chemotherapy. Other drugs, notably chloramphenicol and benzol, seem to operate in an idiosyncratic fashion; many other agents have been implicated, though each in fewer cases. Defective synthesis of deoxyribonucleic acid (DNA) and abnormal nuclear maturation result from malabsorption of vitamin B12 (as in pernicious anemia) or dietary deficiency of folic acid or its malabsorption (sprue). Methotrexate, fluorouracil, hydroxyurea, mercaptopurine, and other drugs used in chemotherapy affect purine and pyrimidine metabolism and can produce the same lesion. All of these anemias are generally characterized by large red blood cells, which are called macrocytes, and a specific morphological abnormality of the nuclear chromatin of erythroblasts which characterizes them as megaloblasts, and the condition as megaloblastic anemia. See VITAMIN B12. Defective synthesis of hemoglobin impairs cytoplasmic maturation. The majority of cases are due to deficiency of body stores of iron and to abnormal release of iron from reticuloendothelial stores. The former occurs in iron-deficiency anemia, and the latter in the anemia of chronic inflammatory diseases. Qualitative abnormalities of globin synthesis occur in the hereditary hemoglobinopathies (as in sickle-cell anemia, hemoglobin-C disease), whereas reductions in the rate of synthesis of either the α- or β-globin chains, in the absence of qualitative abnor-
malities, occurs in the thalassemia syndromes. The hemoglobinopathies and thalassemia may be inherited independently of one another or simultaneously. Heterozygotes have little if any anemia, whereas homozygotes and double heterozygotes have moderate or marked anemia. Rarely, disorders affect the porphyrin moiety of the heme prosthetic group in globin (as in lead intoxication and hereditary erythropoietic porphyria). Acute blood loss reduces the total blood volume and produces symptoms of weakness, dizziness, thirst, faintness, and shock, in that order, according to increasing magnitude of blood loss. The anemia which results is not detectable by measures of concentration until hemodilution occurs over subsequent days, or more rapidly if replacement fluids are given intravenously. The proliferative response of the healthy marrow will correct the anemia in 2– 6 weeks, depending upon the size of the deficit and provided there are sufficient body stores of iron, folic acid, and vitamin B12, required for hemoglobin synthesis. Chronic blood loss results in iron-deficiency anemia. Hemolysis, the accelerated destruction of red blood cells, also induces a proliferative response from the marrow. However, it differs from hemorrhage because red blood cells are lost without plasma, and it thus diminishes the measures of concentration at the outset. Disorders external to the red blood cell (such as autoantibodies, chemicals including drugs, splenomegaly, irradiation, and thermal injury) are almost always acquired, and may destroy otherwise healthy cells of the host, as well as transfused, normal red blood cells. Even trauma to the red blood cells in the circulation (such as from artificial heart valves, abnormal capillaries, sometimes marching or running on hard surfaces) has caused hemolytic anemia. Disorders intrinsic to the red blood cell are almost always hereditary, and may affect the membrane, or envelope itself (as in hereditary spherocytosis and elliptocytosis), or the glycolytic enzymes of the cytosol (as in hereditary deficiency of pyruvate kinase or glucose-6phosphate dehydrogenase), or be the consequence of anomalies of globin synthesis. When an external factor interacts with the intrinsic defect (such as the spleen: hereditary spherocytosis; a redox drug or chemical: abnormality of glucose-6-phosphate dehydrogenase), the anemia is abated when the external factor is removed, even though the hereditary factor remains. In hemolytic anemia the iron is recovered and used again for hemoglobin synthesis, except when there is intravascular hemolysis, in which case it may be lost in the urine. Diagnosis and treatment. Pallor, weakness, and fatigue are common to all anemias. They may not be noticed until anemia is advanced, if it is of gradual onset and there has been time for cardiovascular and biochemical adaptation. Faint jaundice in the sclerae is a feature of hemolytic anemia, whereas in anemia due to lack of vitamin B12, glossitis and neuropathy occur and may be noted. The morphological study of the blood is often the most simple
Anesthesia guide. Macrocytic anemias are often found to be megaloblastic and due to deficiency of vitamin B12 or folic acid. The latter contention can be proved by measuring the quantity of those substances in the plasma. Administration of one of those vitamins will only cure individuals in whom its specific deficiency is established. Microcytic anemias are often also hypochromic. Hypochromia is a reduction in the concentration of hemoglobin within individual red blood cells. It is detected by microscopic examination of the thin blood film after Wright’s stain, and by finding the ratio of hemoglobin concentration/hematocrit to be less than 0.31. Microcytic-hypochromic anemias are most often due to iron deficiency. The administration of iron will resolve iron-deficiency anemia. Since the body normally conserves iron and has no excretory route for it, iron deficiency implies a pathological blood loss which must be explained. The administration of iron or iron-containing over-the-counter medications without investigation of the underlying cause may result in the failure to detect an asymptomatic neoplasm of the gastrointentinal tract which is causing occult blood loss until it has progressed to the incurable state. In the absence of proved deficiency, the administration of iron, vitamin B12, and folic acid are not of value to the anemic individual. Further, the administration of iron, especially by injection, may be harmful. Prednisone and other adrenal corticosteroids are helpful in hemolytic anemias associated with autoantibodies. Hereditary disorders are generally not amenable to therapy, except for those hemolytic diseases which may benefit after splenectomy. In all other cases, the treatment of the anemia is achieved by treating the underlying disease, such as hypothyroidism, rheumatoid arthritis, or leukemia. Blood transfusions are reserved for acute blood loss when symptoms of hypovolemia and shock are present, or in chronic anemia if there are signs of inadequate cardiovascular or pulmonary compensation and an underlying cause cannot be found or treated. Since blood transfusions entail a risk of immune reactions, hepatitis, acquired immune deficiency syndrome, and iron overload, they should be avoided when other alternatives for treatment are possible. See CLINICAL PATHOLOGY; HEMATOLOGIC DISORDERS. Arthur Haut Bibliography. G. R. Lee et al., Clinical Hematology, 8th ed., 1981; W. J. Williams et al. (eds.), Hematology, 3d ed., 1983.
Anesthesia Loss of sensation with or without loss of consciousness. There are several ways of producing anesthesia, with the choice dependent on the type of surgery and the medical condition and preference of the patient. Each person responds differently to a given anesthetic, and anesthetic techniques and drugs often have marked effects on bodily functions, especially those of the cardiovascular and respiratory systems. Therefore, these systems are moni-
tored closely during anesthetic administration, with measurements such as heart sounds, blood pressure, heart action (electrocardiogram), temperature, and oxygenation taken using a variety of sophisticated devices. General anesthesia. A state of complete insensitivity or unconsciousness is produced when anesthetic gases are inhaled; adjuvant drugs are often given intravenously. Anesthetics commonly used by inhalation are halogenated-hydrocarbon liquids (isoflurane, enflurane, and halothane), which are vaporized and administered in oxygen, usually with nitrous oxide. Adjuvant drugs include thiopental (to produce rapid loss of consciousness for induction of anesthesia), narcotic opioids (for supplemental analgesia, that is, insensitivity to pain), muscle relaxants, and sedatives. Although the mechanism of general anesthesia is unknown, the anesthetics act on the upper reticular formation of neurons in the thalamus and midbrain (neuronal structures necessary for activating the cerebral cortex and maintaining an active, attentive state). Regional anesthesia. Analgesia, without loss of consciousness, results from injecting a solution of local anesthetic drug either into the cerebrospinal fluid surrounding the spinal cord (spinal anesthesia) or into the epidural space surrounding the cerebrospinal fluid (epidural anesthesia), usually in the lower back (lumbar) area. Epidural injection is also performed at the base of the spine (caudal anesthesia). The duration and extent of analgesia depend on the nature and amount of local anesthetic injected, among other factors. The local anesthetic acts by blocking the conduction of nerve impulses. Narcotic opioids are injected postoperatively into either the epidural space or cerebrospinal fluid for pain relief. Local anesthesia. Analgesia can be localized to a small area, for example, the forearm, by injecting a local anesthetic solution around nerves supplying the area (the bracheal plexus in the upper arm and chest supplies the forearm). Large peripheral nerves may also be blocked individually by using this method. The duration of analgesia is determined by the choice and amount of local anesthetic. Injection of a dilute solution of local anesthetic into the skin and superficial tissues provides very localized analgesia for minor surgery. Acupuncture. Acupuncture is an ancient procedure, once used only in China but now practiced in the United States and elsewhere. It involves inserting needles into specific points around the body, as determined from historical charts, and manipulation of the needles; sometimes electric current is applied. Weak analgesia results through alteration of pain perception. Although sometimes helpful for chronic pain, acupuncture generally has not been found satisfactory for surgical anesthesia. See CENTRAL NERVOUS SYSTEM; PAIN. Fredrick K. Orkin Bibliography. M. J. Cousins and P. O. Bridenbaugh (eds.), Neural Blockade in Clinical Anesthesia and Management of Pain, 3d ed., 1998; R. D. Dripps, J. E. Eckenhoff, and L. D. Vandam (eds.), Introduction to Anesthesia: The Principles of Safe Practice,
665
666
Aneurysm 7th ed., 1988; R. D. Miller (ed.), Anesthesia, 4th ed., 1994; R. K. Stoelting, Pharmacology and Physiology in Anesthetic Practice, 3d ed., 1998.
Aneurysm A localized, abnormal arterial dilation usually caused by a weakening of the vessel wall. Aneurysms are commonly found in the abdominal aorta, intracranial arteries, and thoracic aorta; however, they may also involve the femoral, popliteal, splenic, carotid, or renal arteries. Aneurysms vary in size from less than 1 in. (2–3 cm) to more than 4 in. (10 cm). They are usually classified as true or false: true aneurysms involve all layers of the artery (inner endothelium or intima, middle muscular layer or media, and outer connective tissue or adventitia); false aneurysms do not involve all layers. With the exception of intracranial aneurysms, atherosclerotic vessel disease is generally the cause of aneurysms; other causes include syphilis, trauma, cystic medial necrosis, bacterial infections, and arteritis. Complications. The major clinical sign of all aneurysms (except intracranial) is a large pulsatile mass; other manifestations depend on the location of the aneurysm. Aneurysms of the thoracic aorta can produce symptoms of compression or erosion as well as dissection. In dissection, the blood creates a channel between the intima and media of the aorta, which may result in blockage of major branches of the aorta, a condition that is often fatal. Debris in aneurysms may also embolize; that is, it may break free and block distal arteries. The major complication of an aneurysm is rupture, the possibility of which is directly related to its size. The mortality rate for a ruptured abdominal aortic aneurysm is 60–70%. Intracranial aneurysms, often called berry aneurysms, are congenital weaknesses in the intracranial arteries. These aneurysms are saclike outpocketings of the vessel and vary in size from 0.2 in. (0.4 cm) to greater than 0.6 in. (1.5 cm). The major manifestations are usually related to bleeding and can vary from a severe headache, to neurologic impairment, to death. In many instances, a small initial hemorrhage (a “herald bleed”) is often followed by a second bleeding episode, which is much more severe and often fatal. See HEMORRHAGE. Diagnosis and treatment. The best method of diagnosis is ultrasound or a computer x-ray analysis (computed tomography scan). The recommended treatment for all aneurysms is surgery. In patients whose aneurysms are due to atherosclerosis, associated diseases can increase the risk of surgery. Abdominal and thoracic aneurysms are generally treated with resection of the aneurysm and replacement with a graft made of synthetic material (see illus.). Blood flow must be established to the branches of the aorta. Intracranial aneurysms are also treated surgically; however, surgery involves excluding the aneurysm from the cerebral circulation by placing a metallic clip across the neck of the aneurysm sac. A ruptured aneurysm or dissecting aneurysm is a true
(a)
(b)
Surgical repair of an abdominal aortic aneurysm. (a) Before resection. (b) Tubular-type graft anastomosis. (After M. D. Kerstein, P. V. Moulder, and W. R. Webb, eds., Aneurysms, Williams and Wilkins, 1983)
emergency and requires immediate surgical intervention. See ARTERIOSCLEROSIS; CARDIOVASCULAR SYSTEM; CIRCULATION DISORDERS; HYPERTENSION; MEDICAL IMAGING. Donald L. Akers; Morris D. Kerstein Bibliography. C. Arbost, M. Allary, and N. Economou, Resection of an aneurysm of the abdominal aorta: Reestablishment of the continuity by a preserved human arterial graft which results after five months, Arch. Surg., 64:405, 1952; M. D. Kerstein, P. V. Moulder, and W. R. Webb (eds.), Aneurysms, 1983; S. I. Schwartz and G. T. Shires (eds.), Principles of Surgery, 7th ed., 1998.
Angiogenesis The origin and development of blood vessels. Blood vessels are composed of two basic cell types, vascular endothelial cells and peri-endothelial cells (including vascular smooth muscle cells and elongated contractile cells called pericytes, both of which support the underlying endothelial cells). The inner epithelial lining of all blood vessels, adjacent to the lumen, is a single layer of endothelial cells (Fig. 1). In larger blood vessels, such as arteries and veins, the inner endothelial lining, called the tunica intima, is surrounded by a medial layer, the tunica media, composed of multiple layers of vascular smooth muscle cells embedded in elastin-rich extracellular matrix. The tunica media layer is surrounded by an extracellular matrix-rich layer called the tunica adventitia. In contrast, capillary walls consist of only a single layer of endothelial cells, sometimes surrounded by pericytes. Arteries and veins are the two fundamental types of blood vessels, carrying blood away from and toward the heart, respectively. Although both have the basic structure noted above, higher-pressure arteries generally have thicker medial smooth muscle– containing layers, whereas larger veins have thinner, more elastic walls and valves to prevent backflow. Recent studies of the process of blood vessel formation, or angiogenesis, have shown that the endothelial cells contributing to arteries and veins have
Angiogenesis
fibrous connective tissue (tunica adventitia)
fibrous connective tissue (tunica adventitia)
smooth muscle (tunica media)
smooth muscle (tunica media)
elastic layer
elastic layer
endothelium
endothelium (tunica intima)
endothelial cells
endothelial cells (tunica intima) pericyte (a)
(b)
(c)
Fig. 1. Structure of blood vessels. (a) Arteries and (c) veins have the same tissue layers. (b) Capillaries are composed of only a single layer of endothelial cells, sometimes surrounded by elongated contractile cells known as pericytes (not to scale). (Adapted from P. H. Raven and G. B. Johnson, Biology, 6th ed., McGraw-Hill, New York, 2002 )
distinct molecular and functional identities that are established early during development. Blood vessel formation. Blood vessels are found only within the chordates, but within the vertebrates the anatomical architecture of the major blood vessels of the circulatory system is surprisingly well conserved, as highlighted by recent detailed characterization of the anatomy of the developing zebrafish vasculature. The first major vessels to emerge during embryonic development form by coalescence of individual mesodermal endothelial progenitor cells, or angioblasts, into cords of attached cells which then form open vascular cells, or lumenize. This process, called vasculogenesis, was thought to be restricted to early embryonic development; however, recent evidence has shown that vasculogenesis also occurs later in development and even postnatally. Most later developmental and postnatal blood vessel formation, however, occurs by sprouting and elongation of new vessels from preexisting vessels or remodeling of preexisting vessels, collectively known as angiogenesis. Vessels generally form first as unlined endothelial tubes but generally become rapidly associated with supporting pericyte or vascular smooth muscle cells. A variety of studies have shown that the acquisition of these supporting cells is critical for proper morphogenesis, stability, and survival of nascent blood vessels. Although most larger blood vessels form during development and are relatively stable thereafter, significant growth and remodeling of blood vessels continues throughout adult life, for example during uterine cycling or in conjunction with gain or loss of fat or muscle mass. Postnatal vessel growth also occurs in pathologic contexts, notably cancer. Tumors recruit new blood vessels to provide themselves with access to the host circulation and obtain oxygen and nutrients. Acquiring a vascular blood supply is essential for the progression and survival of a tumor. New tumor vessels form primarily through endothelial proliferation and growth from preexisting, adjacent host blood ves-
667
sels. These vessels are usually morphologically abnormal, poorly structured, and deficient in vascular smooth muscle cells. Some tumor vascular spaces appear to lack even an endothelial lining and are surrounded only by tumor-derived cells. Molecular regulators. Recently, interest in factors that might promote or inhibit blood vessel growth and the application of molecular methods to vascular biology and the study of vascular development have led to the discovery of many new genes encoding proteins involved in blood vessel growth and assembly. The vascular endothelial growth factor (VEGF) family of ligands and their tyrosine kinase vascular endothelial growth factor receptors (VEGFRs) expressed on the surface of endothelial cells play critical roles in the differentiation of blood vessels. VEGF-A plays a particularly central role—it is essential for the survival, proliferation, migration, and arterial differentiation of endothelial cells. Targeted disruption or “knockout” of just one of the two copies of the gene encoding this protein in mice causes (heterozygous) embryos to die early in development with a dramatic reduction in endothelial cell number. VEGFR2 is a key endothelial receptor transducing VEGF-A signals, and loss of the gene encoding this protein is also lethal (although only in homozygotes with both functional copies of the gene knocked out). Other related VEGF ligands and VEGFR receptors have also been uncovered, and these also play important roles in modulating the development of blood vessels. VEGF-C signaling through the VEGFR3 receptor plays an essential role in the formation of lymphatic vessels. Another set of ligands, the angiopoietins, play a critical role in vascular endothelial– vascular smooth muscle cell interaction and blood vessel remodeling. Angiopoietin-1 binds to the endothelially expressed Tie-2 receptor to promote the acquisition of vascular smooth muscle cells and stabilization and maturation of nascent blood vessels. Loss of either angiopoietin-1 or Tie-2 function is lethal
Angle modulation Zebrafish adults are small; fully grown animals are only about an inch long, so large numbers of fish can be housed in a small space. Zebrafish have a relatively short generation time, and adult females lay hundreds of eggs every few weeks, so large numbers of progeny can be obtained for genetic or experimental studies. The eggs are fertilized and develop externally to the mother; thus zebrafish embryos and larvae are readily accessible for noninvasive observation or experimental manipulation at all stages of their development. Early development is very rapid: Blood circulation begins within 24 hours after fertilization, larvae hatch by approximately 2.5 days, and larvae are swimming and feeding by 5–6 days. Two additional features of zebrafish make them particularly useful for studying vascular development. First, developing zebrafish are very small— a 2-day postfertilization (dpf) embryo is just 2 mm long. The embryos are so small, in fact, that the cells and tissues of the zebrafish receive enough oxygen by passive diffusion to survive and develop in a reasonably normal fashion for the first 3 to 4 days, even in the complete absence of blood circulation. This makes it fairly straightforward to assess the cardiovascular specificity of genetic or experimental defects that affect the circulation. Second, zebrafish embryos and early larvae are virtually transparent. The embryos of zebrafish (and many other teleosts) are telolecithal; that is, yolk is contained in a single large cell separate from the embryo proper. The absence of obscuring yolk proteins gives embryos and larvae a high degree of optical clarity. Genetic variants deficient in pigment cells or pigment formation are even more transparent. This remarkable transparency is probably the most valuable feature of the fish for studying blood vessels in vivo, and has been exploited in a number of different ways. Confocal microangiography, a technique that permits high-resolution three-dimensional visualization of fluorescent structures, has been used to study functioning blood vessels throughout the zebrafish at every stage of embryonic or early larval development. Transgenic zebrafish expressing green fluorescent proteins in blood vessels allow visualization of vascular endothelial cells and their angioblast progenitors prior to initiation of circulation through a vessel, and even before vessel assembly. A complete atlas of the anatomy of the developing vasculature throughout early development has been prepared using confocal microangiography, the first such atlas to be compiled for any vertebrate. Together these and other tools make it possible to perform in-vivo genetic and experimental dissection of the mechanisms of blood vessel formation in the zebrafish. This facility has already been used to dissect the nature of molecular signals guiding arterial-venous differentiation of endothelial cells during development, and promises to make the fish an important resource for understanding the cues and signals that guide the patterning of growing blood vessels. See ARTERY; BLOOD VESSELS; CANCER (MEDICINE); CAPILLARY (ANATOMY); CIRCULATORY SYSTEM; DEVELOPMENTAL BIOLOGY; TUMOR; VEIN. Brant M. Weinstein
Bibliography. E. M. Conway, D. Collen, and P. Carmeliet, Molecular mechanisms of blood vessel growth, Cardiovasc. Res., 49:507–521, 2001; N. Ferrara, VEGF and the quest for tumour angiogenesis factors, Nat. Rev. Cancer, 2:795–803, 2002; J. M. Isner, Myocardial gene therapy, Nature, 415:234– 239, 2002; S. Isogai, M. Horiguchi, and B. M. Weinstein, The vascular anatomy of the developing zebrafish: An atlas of embryonic and early larval development, Dev. Biol., 230:278–301, 2001; R. Kerbel and J. Folkman, Clinical translation of angiogenesis inhibitors, Nat. Rev. Cancer, 2:727–739, 2002; N. D. Lawson and B. M. Weinstein, Arteries and veins: Making a difference with zebrafish, Nat. Rev. Genet., 3:674–682, 2002; B. Weinstein, Vascular cell biology in vivo: A new piscine paradigm?, Trends Cell Biol., 12:439–445, 2002.
Angle modulation Modulation in which the instantaneous angle varies in proportion to the modulating waveform. When used in an analog manner, it is referred to as either frequency modulation (FM) or phase modulation (PM), depending upon whether the instantaneous frequency or instantaneous phase varies with the modulation. When used in a digital modulation format, it is referred to as either frequency-shift keying (FSK) or phase-shift keying (PSK). Both FSK and PSK can be used with either a binary or an M-ary alphabet, where M is an integer greater than 2. Indeed, the most common forms of PSK correspond to an input alphabet size of either two or four; while for FSK alphabets of two, four, or eight, symbols are typically employed. The spectrum of an angle-modulated signal for an arbitrary modulating waveform (that is, an arbitrary message) is difficult to determine. However, if that modulating waveform is a sine wave of frequency fm, modulating a carrier of frequency fc, then it is straightforward to show that the modulated signal has a discrete frequency spectrum with components at fc, fc ± fm, fc ± 2fm, fc ± 3fm, . . . (illus. a). The amplitude of the nth harmonic of the spectrum, at frequency fc + nfm, is given by the magnitude of Jn(β), a Bessel function of order n and argument β. This quantity β is known as the modulation index. For phase modulation, β is the maximum deviation of the signal phase from the carrier phase, in radians; while for frequency modulation, it is the ratio f/fm, where f is the maximum deviation of the signal frequency from the carrier frequency. See MODULATION. For a PSK waveform modulated by a message that is modeled as a sequence of random binary data, the power spectral density is a continuous function of frequency, f, which has a peak centered at the carrier frequency, fc (illus. b). It has the form A sin2 x/x2, where A is a constant and x is proportional to f−fc. If, instead of phase-shift-modulating a carrier, the message frequency-shift modulates the signal between the carrier frequencies fc1 and fc2, the power spectral
669
Anglesite
amplitude
2 ∆f
fc − 8f m
fc bandwidth frequency
fc + 8fm
(a)
power spectral density
(a)
5 cm
fc (b)
frequency (b)
Anglesite. (a) Crystals on galena from Phoenixville, Pennsylvania. (Bryn Mawr College specimen). (b) Crystal habits (after L. G. Berry and B. Mason, Mineralogy, W. H. Freeman, 1959)
power spectral density
670
fc1
fc 2
frequency (c) Spectral diagrams of angle-modulated signals. (a) Amplitude-frequency spectrum of a signal that is frequency-modulated or phase-modulated by a sine wave, with modulation index β = 5 (after M. Schwartz, Information Transmission, Modulation, and Noise, 4th ed., McGraw-Hill, 1990). (b) Power spectral densities of PSK signal and (c) FSK signal (after B. P. Lathi, Modern Digital Communications Systems, 3d ed., Oxford University Press, 1998).
density has peaks at both these frequencies (illus. c). See FREQUENCY MODULATION; MODULATION; PHASE MODULATION. Laurence B. Milstein Bibliography. M. Schwartz, Information Transmission, Modulation and Noise, 4th ed., 1990; H. Taub and D. L. Schilling, Principles of Communication Systems, 2d ed., 1986; R. E. Ziemer and W. H. Tranter, Principles of Communications, 4d ed., 1994.
Anglesite A mineral with the chemical composition PbSO4. Anglesite occurs in white or gray, orthorhombic, tabular or prismatic crystals or compact masses (see illus.). It is a common secondary mineral, usually formed by the oxidation of galena. Fracture is conchoidal and luster is adamantine. Hardness is 2.5– 3 on Mohs scale and specific gravity is 6.38. Anglesite fuses readily in a candle flame. It is soluble with difficulty in nitric acid. The mineral does not
occur in large enough quantity to be mined as an ore of lead, and is therefore of no particular commercial value. Fine exceptional crystals of anglesite have been found throughout the world. In the United States good crystals of anglesite have been found at the Wheatley Mine, Phoenixville, Chester County, Pennsylvania, and in the Coeur d’Alene district of Shoshone County, Idaho. Edward C. T. Chao
Anguilliformes The true eels, a large order of actinopterygian fishes, also known as the Apodes. The Anguilliformes are related to Saccopharyngiformes (sack pharynx fishes), Elopiformes (tarpons), and Notacanthiformes (spiny eels and halosaurs); all have a leptocephalous larval stage in development. The chief characteristics of the Anguilliformes include an elongate body with numerous vertebrae; a pectoral girdle which, when present, is free from the head; absence of a pelvic girdle and pelvic fins in Recent adults; dorsal and anal fins confluent with the caudal fin (caudal fin absent in some ophichthids); loss of skeletal parts, especially about the head (for example, the orbitosphenoid, posttemporal, symplectic, posttemporal bones, and gular plate are absent); mesocoracoid and postcleithra absent; swim bladder present, usually physostomous (that is, with a connection to the esophagus); restricted gill openings; no fin spines; scales usually absent or, if present, cycloid and embedded in the skin; and 6 to 49 branchiostegal rays. See ACTINOPTERYGII; EEL; ELOPIFORMES; NOTACANTHOIDEI; SACCOPHARYNGIFORMES; SWIM BLADDER. The Anguilliformes comprise three suborders, 15 families, 141 genera, and about 791 species. The largest family of each suborder is treated below.
Angular correlations Suborder Anguilloidei. Frontal bones are divided (sutured), and scales present or absent. There are three families, five genera, and about 29 species. The largest family is the Anguillidae (freshwater eels). Unlike most eels of the world, anguillids have small scales embedded in the skin; however, scales do not develop until an individual is 2 or 3 years old. They are further distinguished by having pectoral fins, and a caudal fin continuous with the dorsal and anal fins; branchiostegals do not overlap; the anterior nostril is tubular, near the tip of the snout; the posterior nostril is round and in front of the eye. One genus and about 16 species occur worldwide, except in mainland Asia, the eastern Pacific, and most of South America. Adult anguillids live in freshwater streams or in brackish-water estuaries and, upon reaching maturity, migrate to sea to spawn.
American eel (Anguilla rostrata). (After G. B. Goode, Fishery Industries of the U.S., 1884)
In the case of the common American eel (Anguilla rostrata) [see illustration], the life cycle encompasses several morphological and physiological stages: egg > leptocephalus > glass eel > elver > yellow eel > silver eel > death after spawning. The act of spawning has never been observed. The leptocephalus is a planktonic larva; the glass eel is the first stage of the metamorphosed leptocephalus and lacks pigment; the elver is completely metamorphosed and pigmented; the yellow eel is the nonmigratory juvenile; and the mature silver eel is ready for seaward migration. The leptocephali drift toward the coast of the Americas, a journey requiring about one year, at which time they attain 60–80 mm (2.4– 3.2 in.) in length and begin to metamorphose. During metamorphosis, larvae shrink, leaving the glass eels shorter than their precursors. Female eels ascend Atlantic and Gulf coastal rivers, often far into the heartland of North America. There they remain 5–40 years (depending on latitude) before returning to sea. In the meantime, males remain on or near the coast. Males are generally smaller than females and develop larger eyes just before migration. It remains a mystery as to the exact place of spawning. Adult American eels have never been captured beyond the continental shelf, and neither spawning adults nor their eggs have been observed. The only empirical evidence of a possible spawning site is in 2000 m (6560 ft) of water of a geologically old, deep basin of the Bahamas where two adults were photographed. Suborder Muraenoidei. Frontals are divided (sutured), and scales absent. There are three families, 24 genera, and about 207 species.
The largest family is the Muraenidae (moray eels). Members of the family, consisting of 15 genera and about 185 species, lack pectoral fins; the gill opening is restricted to a small oval aperture; the teeth are long and fanglike; posterior nostrils are above and in front of the eyes; and some species of Gymnothorax are ciguatoxic (ciguatoxin is a poisonous substance accumulated up the food chain in the flesh and viscera of some fish). Maximum length is about 3 m (9.8 ft). Moray eels occur in the tropical and temperate seas of the world. Suborder Congroidei. Frontals are fused scales rarely present. There are nine families, 112 genera, and about 555 species. The largest family is the Ophichthidae (snake eels and worm eels). Ophichthidae is a family of marine eels, none of which regularly occurs in freshwater. The family differs from other eels in the unusual form of the branchiostegal rays, which are long, slender, and numerous, do not attach basally to the hyoid bone, and broadly overlap those of the other side, leaving the branchiostegal region inflated or slightly bulbous. It is the most diverse and speciose family of eels, consisting of about 250 known species in 55 genera. Ophichthids occur in coastal marine habitats of tropical and warm temperate waters of the world. Herbert Boschung Bibliography. E. B. Böhlke (ed.), Fishes of the Western North Atlantic, Part 9, vol. 1: Orders Anguilliformes and Saccopharyngiformes (pp. 1–655), vol. 2: Leptocephali (pp. 657–1055), Sears Foundation for Marine Research Memoir (Yale University), New Haven, 1989; J. E. McCosker and R. H. Rosenblatt, A revision of the Eastern Pacific snakeeel genus Ophichthus (Anguilliformes: Ophichthidae) with the description of six new species, Proc. Calif. Acad. Sci., 50(19):397–432, 1998; J. S. Nelson, Fishes of the World, 4th ed., Wiley, New York, 2006.
Angular correlations An experimental technique that involves measuring the manner in which the likelihood of occurrence (or intensity or cross section) of a particular decay or collision process depends on the directions of two or more radiations associated with the process. Traditionally, these radiations are emissions from the decay or collision process. However, a variant on this technique in which the angular correlations are between an incident and emitted beam of radiation has been widely used; this variant is known as angular distributions. The fundamental reason for performing such measurements, rather than just scrutinizing a single radiation in a particular direction or measuring the total intensity for a process, is that the angular correlation or angular distribution measurement provides much more information on both the decay or collision process and on the structure and properties of the emitter of the radiation. The technique is used to study a variety of decay and collision processes in atomic and molecular physics, condensed-matter
671
672
Angular correlations (solid-state) and surface physics, and nuclear and particle physics. Detection of multiple emissions. The principal use of this technique in nuclear physics has been to determine the angular momentum, or spin, and parity of excited nuclear states which are radioactive, that is, decay spontaneously, by measuring in coincidence the radiation in specific directions from two successive transitions in the radioactive cascade. The measurements are generally of coincidences between gamma rays, but coincidences between gamma rays and electrons (beta particles) are also used. The form of the angular correlation, the measured intensity as a function of the angle between the two radiations, gives the information about the intermediate excited state in the cascade. See GAMMA RAYS; NUCLEAR SPECTRA; NUCLEAR STRUCTURE; RADIOACTIVITY. In atomic and molecular collisions as well as in nuclear and particle collisions, this technique is employed as a means of completely specifying the dynamics of the collision, with the added proviso that the energies of the emitted radiations are also to be measured. Wide use has been made of angular correlations in the impact ionization of atoms by electrons where the directions of both the scattered electron and the ejected electron are measured. This application is known as (e,2e) spectroscopy. These studies, particularly just above the ionization threshold energy, have revealed the details of the law governing threshold ionization by charged particles, known as the Wannier law. Furthermore, these measurements reveal considerable information about the quantum-mechanical wave function of the system after the ionization process has taken place, that is, the final state. This type of final state, two electrons and one positive ion, is an example of a three-body continuum Coulomb state. See ATOMIC STRUCTURE AND SPECTRA; SCATTERING EXPERIMENTS (ATOMS AND MOLECULES). Angular distributions. This technique is much simpler experimentally than is the general angular correlations technique, since coincidence measurements are not required. Consequently, it can be employed in situations where the probability (or cross section) for the process is so small that a coincidence experiment would not lead to enough counts to make the result meaningful. This technique allows the extraction of the multipolarity of the incident radiation, provided the target orientation is known. The widest application of this technique has been in atomic, molecular, and condensed-matter investigations, and both emitted electrons and electromagnetic radiation (x-rays, ultraviolet, and visible) have been studied. The most important incident particles are electromagnetic radiation, generated principally by synchrotron radiation. Other incident particles include electrons, protons, and heavy ions and atoms. See MULTIPOLE RADIATION; SYNCHROTRON RADIATION. Photoelectron spectroscopy. In photoionization, the study of the angular distribution of photoelectrons emitted from a target as a result of ionization by electromagnetic radiation (photons), the technique is used in several different ways. The atomic and
molecular targets are generally randomly oriented. If unpolarized photons are used, the photoelectron direction with respect to the photon beam direction is measured; for linearly polarized photons, the photoelectron direction with respect to the polarization is measured. This technique is used primarily at photon energies low enough (below about 500 electronvolts, corresponding to photon wavelengths longer than about 2.5 nanometers) that they constitute electric dipole radiation, to an excellent approximation. The general form of the angular distribution is determined by a parameter β, known as the asymmetry parameter, that contains dynamical information about both the target and the photoionization process. The dependence of β on photon energy is extremely helpful in understanding the details of the photoionization cross section, particularly in energy regions where resonances or minima occur. In addition, since β depends on the quantum-mechanical phases of each of the possible final states, phase information, which is unavailable from total cross section measurements, can be obtained. See ELECTRON SPECTROSCOPY; PHOTOIONIZATION; POLARIZED LIGHT. For oriented targets, that is, target atoms or molecules whose magnetic moments, spins, or angular momenta are all aligned, the general form of the angular distribution is more complex and depends upon both the degree of alignment and the state of polarization of the incident radiation. This is often the case for studies of surfaces and crystalline solids where the crystal structure aligns the constituent atoms and molecules; then the technique is generally known as angle-resolved photoelectron spectroscopy. In surface studies, this technique can often aid in the understanding of the details of how the atoms and molecules are attached to the surface along with the geometry of their orientation with respect to the surface. See ANGULAR MOMENTUM; SPIN (QUANTUM MECHANICS); SURFACE PHYSICS. Decay of excited states. When an excited state decays, energy is given off in the form of some kind of radiation. Measuring the angular distribution of the radiation provides certain information about the excited state. In particular, if the emission is isotropic, that is, independent of direction, the excited state was also isotropic. An angular distribution that does depend upon direction (anisotropic) implies that the excited state too had a preferred direction. Thus, the technique provides information about the process that created that excited state. The technique is used in nuclear physics with observation of emitted gamma rays, particularly in connection with radioactive decay. It is also used extensively in atomic and molecular physics. When electromagnetic radiation emitted from atoms or molecules is observed, the technique is called fluorescence spectroscopy, and with emitted electrons it is known as Auger spectroscopy. See AUGER EFFECT; X-RAY FLUORESCENCE ANALYSIS. Steven T. Manson Bibliography. J. Boyle and M. Pindola (eds.), Many-Body Interactions in Atomic Physics, 1995; P. C. Deshmukh, V. Radojeviˇc, and S. T. Manson, Photoionization of the outer shells of radon and radium, Phys. Rev. A, 45:6339–6348, 1992.
Angular momentum according to Eqs. (4). Thus, π · L → + L, so that
Angular momentum In classical physics, the moment of linear momentum about an axis. A point particle with mass m and velocity v has linear momentum p = mv. Let r be an instantaneous position vector that locates the particle from an origin on a specified axis. The angular momentum L can be written as the vector cross-product in Eq. (1) The illustration shows the geometrical meaning of this equation. L = r× p
(1)
See CALCULUS OF VECTORS; MOMENTUM. The time rate of change of the angular momentum is equal to the torque N. A rigid body satisfies two independent equations of motion (the dynamical equations) given by Eqs. (2) and (3), where d/dt d p=F dt
(2)
d L=N dt
(3)
denotes the rate of change, the derivative with respect to time t. Only Eq. (1) is required for a point particle. Equation (2) indicates that a rigid body acts as a point particle located at its center of mass. The motion of the center of mass depends upon the net force F, which is the vector sum of all applied forces. Equation (3) gives the angular motion about the center of mass. The case of statics occurs when the net force and net torque both vanish. See KINETICS (CLASSICAL MECHANICS); STATICS; TORQUE. Vectors and pseudovectors. Given the definition of angular momentum in Eq. (1), although it has three components, it is not a vector. The parity transformation π transforms vectors such as r and p
r = π · r = −r
(4a)
p = π · p = −p
(4b)
L is called a pseudovector or axial vector. Mathematically, these quantities are called tensors of odd relative weight. Examples of pseudovectors include angular momentum, torque, and magnetic fields. Two or more pseudovectors are combined by adding their components, so that the total angular momentum is the sum of the angular momenta of the individual particles. Among the fundamental particles, the photon field is a massless vector field and the rho meson is a pseudovector field. See ELEMENTARY PARTICLE; PARITY (QUANTUM MECHANICS). Rigid-body motion. To study the nonrelativistic mechanics of a rigid body made up of N particles, it is convenient to write the instantaneous velocity vα of the α-th particle in terms of angular velocity ω as vα = ω × rα. For this system, Eq. (1) becomes Eq. (5), L=
N
mα rα × (ω × rα ) = I · ω
(5)
α=1
where I is the moment of inertia tensor. The moment of inertia tensor can be written in matrix form as Eq. (6). Representative components of this tensor are given by Eqs. (7). Since I is symmetric, it can alIij =
N α=1
mα [rα2 δij − (rα )i (rα ) j]
I11 =
N α=1
I12 = −
mα (y2α + zα2 ) N
mα xα yα
(6)
(7a)
(7b)
α
angular momentum (L) body
R (vector r) θ velocity (v)
reference point Geometrical definition of angular momentum. If we choose a reference point and draw a line R from it to a moving body, then the line R and the velocity v of the body of mass m define a plane. The angular momentum is a vector that lies perpendicular to this plane. If the momentum of the body is p = mv and the line R is a vector r, then the magnitude of the angular momentum L is pr sin θ.
ways be diagonalized and expressed in principal axis (or symmetry axis) form, IPA = diag (I1, I2, I3), and in this form it is constant in time: I˙PA = 0. The Euler angle coordinates (α, β, γ ) are useful because they can be used to rotate a rigid body into its principal axis frame. See EULER ANGLES; MOMENT OF INERTIA; ROTATIONAL MOTION. When all of the forces that act on a system are conservative (that is, ∇ × F ≡ O), a potential energy function V(r) exists for which F = −∇V(r). Consider a rigid body of mass M and inertia tensor I subject to two conservative forces Fi = −∇ RV1(R) and F2 = −∇ α0V2(α 0), where R is the position vector of the center of mass, and α 0 = (α, β, γ ) are moment-armweighted Euler angles. The quantities ∇ R and ∇ α0 are gradient operators. The equations of motion (2) and (3) become Eqs. (8) and (9). Using IPA and writMR = F1
(8)
I · ω˙ = N2 = r × F2
(9)
ing out Eq. (9) in component form in the rotating principal-axis frame yields Eqs. (10).
673
674
Angular momentum I1 ω ˙ x − (I2 − I3 )ωy ωz = Nx
(10a)
I2 ω ˙ y − (I3 − I1 )ωx ωz = Ny
(10b)
˙ z + (I1 − I2 )ωx ωy = Nz I3 ω
(10c)
See ENERGY; RIGID-BODY DYNAMICS. Conservation. A symmetry is a transformation that leaves a physical system unchanged. A physical quantity is called invariant under a transformation if it remains the same after being transformed. For example, the solutions to Eqs. (2) and (3), or equivalently Eqs. (8) and (9), are invariant under change of the coordinate origin or orientation of the i, j, and k axes. The freedom to choose any orientation of coordinate axis is called rotational invariance, because one choice of axes can be rotated into another. In physics, the rotational invariance follows from the isotropy and homogeneity of space that has been experimentally established to high accuracy. The diagonalization transformation I → I PA is carried out by rotating the axes, and thus is a consequence of rotational invariance. The principal axes of a rigid body are called the symmetry axes because their mechanical description is simplest there. (Rotational invariance says that the mechanics can be solved in any orientation of the axes but not that an arbitrary orientation is simplest.) The study of symmetry shows that one of the deepest relations in physics is that between dynamics and conservation. A physical quantity B is conserved if Eq. (11) is satisfied; that is, B is constant in time dB =0 dt
(11)
although it may vary in space. Noether’s theorem states that if a physical system is invariant under a continuous symmetry, a conservation law exists, provided that the observable in question decreases rapidly enough at infinity. Thus, when the force is zero everywhere (the system is invariant under translation in space), the linear momentum is conserved. If the torque is zero everywhere (the system is invariant under rotation), the angular momentum is conserved. If the system is invariant under translations in time, the total energy is conserved. The converse to Noether’s theorem is false, as conserved quantities do not imply continuous symmetries. For example, the parity transformation π of Eqs. (4) can be conserved but is not continuous. The set of all three-dimensional rotations form a group called SO(3): S for special (det = +1), O for orthogonal (length-preserving), and 3 for the space dimension. The set of 2 × 2 complex, unitary (lengthpreserving) special transformations is called SU(2). There is a 2:1 homorphism of SU(2) → SO(3), and SU(2) is called the covering group of SO(3). Angular momentum is important in the evolution of celestial objects. The shapes of these bodies and collections of these bodies, spiral galaxies for instance, follow from the space-time variation of their torques. There is a puzzling problem along these lines: If the universe can rotate, it is not clear why it rotates so slowly; that is, why the anisotropics are so
small. A proposed solution is based on the hypothesis of an inflationary epoch in the very early universe, about 10−35 s after the big bang. See BIG BANG THEORY; CONSERVATION OF ENERGY; CONSERVATION OF MOMENTUM; COSMOLOGY; SYMMETRY LAWS (PHYSICS); UNIVERSE. Quantum angular momentum. Quantum mechanics has a richer and more complicated structure than classical physics. Because of this, the relationship between symmetry and conservation is even more useful. Whereas in classical physics the observables and states coincide, in quantum mechanics the states of the system correspond to probability amplitudes ψ and the observables A to self-adjoint operators (hermitian operators with suitable boundary conditions). According to the Heisenberg uncertainty principle, all observables cannot be simultaneously known, rather only those that commute. Two commuting self-adjoint operators can be diagonalized by a single unitary transformation. The complete set of commuting observables specifies the maximum amount of sharp information possible about a quantal system. The time dependence of a state, represented in x-space, where x is the position coordinate, as ψ(x, t), is given by the Schr¨ odinger equation (12), where Hˆ is the system hamiltonian ∂ψ ˆ Hψ(x, t) = i (x, t) (12) ∂t (energy) operator and = 1.05 × 10−34 joule-second and is Planck’s constant divided by 2π. See NONRELATIVISTIC QUANTUM THEORY; QUANTUM MECHANICS. Rotational invariance under SO(3) transformations in quantum mechanics implies the angular momentum commutation relations for a single particle given in vector form by Eq. (13). The z component of this equation can be written as Eq. (14), and the x J × J = i J
(13)
[ Jx , Jy ] = i Jz
(14)
and y components are expressed in similar fashion, where [A, B] is the commutator of A and B, defined by Eq. (15). The quantum-mechanical momentum [A, B] = AB − BA
(15)
J is the most general observable satisfying Eq. (13). The orbital angular momentum L = i r × ∇, where ∇ is the gradient operator, is a special case of a quantum-angular momentum. The total angular momentum J = Jxi + Jyj + Jzk has squared length J2 = Jx2 + Jy2 + Jz2. It is straightforward to use the operator identity [A2,B] = A[A,B] + [A,B]A together with Eq. (13) to show that Eq. (16) is satisfied for [ J 2 , Ji ] = 0
(16)
i = x, y, z. An inspection of Eqs. (13) and (16) shows that the angular momentum complete set of commuting observables is { J2, Ji}, where any single component Ji can be simultaneously diagonal with J2. Conventionally Jz is chosen as the diagonal component, and J2 is called the Casimir operator of SO(3).
Anharmonic oscillator The two nonself-adjoint ladder operators J± = Jx ± Jy are useful for determining the properties of angular momentum states. Let ψ( J, M, r) be a simultaneous eigenfunction of the operators J2 and Jz with eigenvalues λ and µ, so that Eqs. (17) and (18) are satisfied. The variable r in J 2 ψ( J, M, r) = λψ( J, M, r)
(17)
Jz ψ( J, M, r) = µψ( J, M, r)
(18)
ψ stands for all other observables which commute with J2 and Jz. Repeated use of the angular momentum commutation relations of the operators J+ and J2 proves that λ = J( J + 1)2 and that µ = M, where M = −J, −J + 1, . . . , J − 1, J for each value of J. Further, J = n/2, where n is a positive integer, so that 0, 1/2 , 3/2 , 1, 5/2 , . . . are the allowed values of total angular momentum. The eigenfunctions for half-oddinteger values of total angular momentum are called spinors. A value S of spin, together with Sz, labels intrinsic spin states. Spinors are new features of rotational invariance. They have the surprising property of changing sign under an SO(3) rotation through 2π radians, R(2π); that is, Eq. (19) is satisfied. The states R(2π)ψ( J, M, r) = −ψ( J, M, r)
expressed by Eq. (21), where j |T( J )| j is the j m|T( JM )| jm = C( j J j , mM) j |T( J )| j (21) reduced matrix element and is independent of mMm. The angular momentum conservation is contained in the Wigner coefficient; the details of the particular operator T are contained only in the reduced matrix element. The operator T in atomic spectroscopy is a perturbation hamiltonian from an incident electromagnetic wave that drives atomic transitions. Thus, the Wigner-Eckart theorem, through the Wigner coefficients, determines which transitions vanish, that is, the selection rules governing the process. Much of atomic spectroscopy follows from the rotational symmetry, or the quantum theory of angular momentum. See ATOMIC STRUCTURE AND SPECTRA; SELECTION RULES (PHYSICS). Brian De Facio; John L. Safko Bibliography. L. C. Biedenharn and J. D. Louck, Angular Momentum in Quantum Physics, 1981; E. U. Condon and H. Odabasi, Atomic Structure, 1980; E. P. Wigner, Group Theory and Its Application to the Quantum Mechanics of Atomic Spectra, 1959.
(19)
ψ( J, M, r) can be spinor-valued since measured probability densities depend only upon |ψ( J, M, r)|2dV (where dV is a volume element), whereas the observables must be tensors to have sensible classical limits. The doublet Zeeman splitting of an unpaired electron in an external magnetic field is an indication that the electron has this strange internal structure. Thus, J = L + S is the total angular momentum. See SPIN (QUANTUM MECHANICS). Quantum addition of angular momenta. The discussion above gives the properties of a single angular momentum. Two angular momenta, J1 and J2, are called independent if all components of one commute with all components of the other. From the preceding discussion, the sets of eigenfunctions ψ( J1, M1, r) and ψ( J2, M2, r) can be determined. Let J = J1 + J2. Then the eigenfunction ψ( J, M, r) can be expressed as a linear combination of the ψ( J1, M1, R) and ψ( J2, M2, r) by Eq. (20), where ψ( J, M, r) = C( J1 J2 J, M1 M2 )ψ( J1 , M1 , r) ·ψ( J2 , M2 , r) (20) C( J1 J2 J, M1M2) is the Wigner coefficient (often erroneously called the Clebsch-Gordon coefficient). The summation is over values of M1 and M2 such that M = M1 + M2, and J must have one of the set of values J = J1 + J2, J1 + J2 − 1. The Wigner coefficients play another important role in the Wigner-Eckart theorem. Let T( JM ) be a family of tensor operators with angular momentum JM, and consider the quantum matrix elements j m |T( JM )|jm. These matrix elements represent transition probability amplitudes from quantum state jm to state jm caused by T( JM ). The Wigner-Eckart theorem states that these matrix elements can be
Anharmonic oscillator An oscillator that does not obey Hooke’s law. This law is an idealized expression that assumes that a system displaced from equilibrium responds with a restoring force whose magnitude is proportional to the displacement. The use of Hooke’s law results in a linear equation of motion that fails to describe many properties of the real world. Nature demonstrates two fundamentally different forms of nonlinearity, which may be called elastic anharmonicity and damping anharmonicity. To understand their difference, and the nature of a harmonic oscillator, it is necessary to understand potential functions. See HARMONIC MOTION; HARMONIC OSCILLATOR; HOOKE’S LAW; NONLINEAR PHYSICS. Atomic interaction. In the classical formulation, where the simplest case is one-dimensional, the xvariable in the potential function V(x) represents displacement away from the stable position of equilibrium that was established between competing forces of attraction and repulsion. For example, in the case of two atoms constrained against rotation, the potential describing their interaction can be approximated by the asymmetric potential curve in Fig. 1a. It is not symmetric because of the Pauli exclusion principle, involving the overlap of electronic orbitals. Thus, it is more difficult to push the atoms together than it is to pull them apart. Therefore, the potential rises more quickly for compression than it does for extension. See EXCLUSION PRINCIPLE; INTERMOLECULAR FORCES; MOLECULAR STRUCTURE AND SPECTRA. Although the potential of Fig. 1 was generated by means of an empirical formula, based on classical physics, V(x) is foundational to the Hamiltonian with which the Schr¨ odinger equation is derived. The Schr¨ odinger equation, in turn, is used to calculate the probability of observing various states of position or
675
Anharmonic oscillator thermal expansion 0
T1
harmonic fit
F (r)
V (r)
T2 >T1 harmonic fit
r0 r0 0.25
0.35
0.45 0.55 r, nm
0.65
0.75
0.32
0.34
0.36
0.38
0.40
0.42
r, nm
(a)
(b)
Fig. 1. Interaction between atoms. (a) Potential function of the interaction, V, plotted against r, the separation distance of the atomic nuclei. V(r) is plotted on a linear scale in units of negative joules. (b) Force, F(r), that is associated with the potential and results from its negative derivative [F(r) = −dV/dr] when the atoms are displaced from the equilibrium position r0. F(r) is plotted on a linear scale in units of newtons.
momentum. Shown superposed on the anharmonic potential of Fig. 1 is a parabola that represents the harmonic approximation. An appreciation for its limitations can be realized by considering two different oscillatory motions, whose turning points are illustrated by the pair of horizontal lines with arrows on each end. These correspond to two different temperatures, T1 and T2, and the amplitude of oscillation is greater at the higher temperature, T2. The shift toward larger time-averaged separation of the atoms as the temperature increases is the basis for thermal expansion. See NONRELATIVISTIC QUANTUM THEORY; ¨ QUANTUM MECHANICS; SCHRODINGER’S WAVE EQUATION; THERMAL EXPANSION. There are numerous other cases in which the harmonic potential is known to be inadequate. To cite one example, its use precludes the proper description of the heat capacity of solids. See SPECIFIC HEAT OF SOLIDS. Simple pendulum. Elastic nonlinearity at large amplitudes of the motion is exhibited also by many macroscopic-scale oscillators, such as the pendulum. Figure 2 shows the potential for a rigid, planar sim-
harmonic fit
unstable equilibrium pendulum potential
stable equilibrium V (θ)
676
−4
−2
0 2 4 angle (θ), radians
6
Fig. 2. Potential function for the rigid, planar simple pendulum.
ple pendulum, in which the bob mass is able to go past vertical, that is, past the unstable equilibrium at the angle θ = π radians. This is not possible for a mass on a string, since the string can support only tension and not compression. See PENDULUM. The pendulum has become a macroscale archetype of chaos. For chaotic motion, the energy with which the damped oscillator is driven must be large enough to allow winding modes, in which the bob is allowed to move “over the top.” Otherwise, it is a system demonstrating that nonlinearity is a necessary but not sufficient condition for chaos. At low levels, but not so small as to introduce complications due to damping nonlinearities, the motion of the pendulum is nearly harmonic and thus isochronous, meaning that the period is independent of amplitude. At high levels, the motion can be periodic and yet highly nonlinear. Highly nonlinear cases involve subharmonic response, in which it takes longer for the pendulum to repeat its motion than is required for the drive. Such modes cannot be explained by a linear equation of motion. See CHAOS. Highly nonlinear but nonchaotic pendulum motion is a realm of amplitude jumps, as drive of the instrument is “tuned” back and forth in frequency across “resonance.” Unlike the harmonic oscillator, when the pendulum is driven at large levels there is not a well-defined resonance frequency. If the drive frequency is below the pendulum’s small-motion resonance frequency, small phase perturbations introduced by an external disturbance can seriously disrupt the steady-state motion. The pendulum can transition from a high-level, long-period mode to a low-level, short-period mode, even though the driving force was not altered. Such a change is suggestive of an analogy in physiology, ventricular fibrillation. See HEART DISORDERS. Damping anharmonicity and internal friction. Anharmonicity manifests itself in different ways at different scales. At both the atomic scale and the macroscopic scale, the assumption of a smooth potential function is oftentimes reasonable. At the mesoscale,
Anharmonic oscillator
V (x)
x
×10 enlargement
x
Fig. 3. Anharmonic potential responsible for internal friction damping. Only at low levels (as noted by the ×10 enlargement) does the nonlinear nature of the potential become observable.
sis, and it involves alterations in the configuration of defect structures. The amount of a given shift is a function of the amplitude of motion in the half cycle that just preceded the shift. For common exponential damping, the shift is proportional to amplitude. An energy-based description of exponential damping is provided by the equation below for the friction force shown in Fig. 5, which involves the function Friction force = (per unit mass) ω v Q πω 2 2 ω x + v2 sgn(v) 4Q
for viscous damping
for hysteretic damping
amplitude change in one cycle
∆ x
sgn(v). [When v > 0, sgn(v) = 1, and when v < 0, sgn(v) = −1.] Here, v is the velocity, ω is the angular frequency, and Q is the quality factor. See CREEP (MATERIALS); Q (ELECTRICITY). Unlike viscous damping which is linear, hysteretic damping is nonlinear. Since the term under the
amplitude change in one cycle
(a)
V (x) metastable equilibria
V (x)
V (x)
this is no longer true, as defect organization takes place via the exceedingly complex interactions responsible for friction. When nature fills the vacuum it abhors, it rarely does so with perfection. At finite temperatures, for example, solids contain lattice sites devoid of an atom. These vacancies are called Schottky defects. Materials are never as strong as theoretical estimates of strength, based on the assumption of perfect lattices, because of defect structures such as dislocations. See CRYSTAL DEFECTS. When material is strained in a cyclic manner, stress/strain hysteresis is observed, and it has been thoroughly studied by engineers. The anharmonic nature of this hysteresis has been described in terms of its influence on oscillator damping as a mechanism of internal friction. It differs in type from elastic anharmonicity, discussed above, and is illustrated by the graph of Fig. 3. See HYSTERESIS. This form of anharmonicity has been the source of challenges to those who seek to study gravitational waves. Metastabilities, illustrated by the metastable equilibria in the enlargement in Fig. 3, are present in the springs used to isolate the mirrors of the LIGO (Laser Interferometer Gravitational-Wave Observatory) instrument from earth noise. Their presence disallows the occurrence of an absolute equilibrium position. A primary problem with trying to explain friction from first principles is that the fine structure responsible for these metastabilities is not static. Rather, because they are both temperature- and stress-dependent, defect structures evolve in time. See GRAVITATIONAL RADIATION; LIGO (LASER INTERFEROMETER GRAVITATIONAL-WAVE OBSERVATORY). Damping anharmonicity is of a totally different character than viscous damping, and only the latter case is well known to physics. Because damping anharmonicity is incompatible with the Hooke’s-law description of a spring, the description of energy loss must be modified (Fig. 4). See DAMPING. In the internal friction case (Fig. 4a), it is seen that the minimum of the potential shifts back and forth as the spring is cycled in strain during free decay. This is due to material creep responsible for hystere-
677
(b)
x
Fig. 4. Comparison of nonlinear and linear oscillator damping. (a) Nonlinear damping due to internal friction. (b) Linear damping due to an external force.
Anhydrite
friction force, hysteretic damping
position (x)
arbitrary linear units
678
time
friction force, viscous damping
velocity (v) Fig. 5. Comparison of (1) viscous (linear) damping that derives from an external-to-the-oscillator friction force and (2) hysteretic (nonlinear) damping of the internal friction type. Both oscillators are in free decay with quality factor Q = 10 and angular frequency ω = 3.07 Hz.
square root in the equation above is twice the energy of oscillation per unit mass, the velocity is seen to be important only by way of its algebraic sign. The two forms are similar only in that they yield the same quality factor Q for the exponential free decay. If one considers a Fourier series to describe the nonlinear friction force, it can be shown that the fundamental component of the series is the only term responsible for the energy loss per cycle that determines the Q. For this reason, many have failed to recognize the ubiquitous nature of nonlinear damping. The two forms differ greatly with regard to frequency dependence of Q. Because ω/Q = 2β, where β is the coefficient commonly referred to as the damping constant, it is seen that Q must be proportional to frequency in the case of viscous damping. When monitored at low frequency and low energy, the vast majority of mechanical oscillators exhibit a Q dependence that is quadratic in the frequency, which is consistent with the nonlinear damping model presented. See FOURIER SERIES AND TRANSFORMS. The properties of internal friction responsible for Q ∝ ω2 were first discovered early in the twentieth century, and the investigators claimed that internal friction in solids appeared to be universal. Experiments since the 1990s support this belief, which implies a closer connection than previously recognized to the well-known case of surface friction described by C. A. Coulomb. In an additional early study by L. Portevin and F. Le Chatelier, which has been also mostly overlooked by physics, it was found that strain of alloys under stress is not smooth but jerky. The discontinuous changes observed are reminiscent of the Barkhausen effect, also discovered early in the twentieth century. See BARKHAUSEN EFFECT; FRICTION. This jerky behavior, related to stick-slip features of friction, implies that atomic-scale quantum mechanics is not best suited to the description of internal friction. The electronvolt is infinitesimal compared to the quantum changes observed. There is evidence for an energy of mesoscale quantization, given by mc2/α, where m = 9.1 × 10−31 kg is the
electron mass, α = 1/137 is the fine-structure constant, and c = 3 × 108 m/s is the speed of light. This Compton scale of energy, having the mesoscale value of 1.1 × 10−11 J, is thought to regulate, through self-organizing metastabilities, the processes “still unknown from first principles” responsible for internal friction. It is possible that one of the last frontiers of physics actually lies in one of the most accessible parts of our world. Randall D. Peters Bibliography. G. L. Baker and J. P. Gollub, Chaotic Dynamics, an Introduction, 2d ed., Cambridge University Press, 1996; R. Desalvo et al., Study of quality factor and hysteresis with the state-of-the-art passive seismic isolation system for Gravitational Wave Interferometric Detectors, Nucl. Instrum. Meth. Phys. Res., Sec. A, 538:526–537, 2005; C. Kittel, Introduction to Solid State Physics, 8th ed., Wiley, 2005; M. R. Matthews, C. F. Gauld, and A. Stinner (eds.), The Pendulum: Scientific, Historical and Educational Perspectives, Springer, 2005; R. D. Peters, Friction at the mesoscale, Contemp. Phys., 45:475–490, 2004; R. D. Peters, Resonance response of a moderately driven rigid planar pendulum, Amer. J. Phys., 64:170–173, 1996.
Anhydrite A mineral with the chemical composition CaSO4. Anhydrite occurs commonly in white and grayish granular masses, rarely in large, orthorhombic crystals (see illus.). Fracture is uneven and luster is pearly to vitreous. Hardness is 3–3.5 on Mohs scale and specific gravity is 2.98. It fuses readily to a white enamel. It is soluble in acids and slightly soluble in water. Anhydrite is an important rock-forming mineral and occurs in association with gypsum, limestone, dolomite, and salt beds. It is deposited directly by evaporation of seawater of high salinity at or above 108◦F (42◦C). Anhydrite can be produced artificially by dehydration of gypsum at about 390◦F (200◦C). Under natural conditions anhydrite
Animal communication
1 cm Anhydrite from Montanzas, Cuba. (Specimen from Department of Geology, Bryn Mawr College)
hydrates slowly, but readily, to gypsum. It is not used as widely as gypsum. Anhydrite is of worldwide distribution. Large deposits occur in the Carlsbad district, Eddy County, New Mexico, and in salt-dome areas in Texas and Louisiana. See GYPSUM; SALINE EVAPORITES. Edward C. T. Chao
Animal Any living organism which possesses certain characteristics that distinguish it from plants is a member of the animal kingdom. There is no single criterion that can be used to distinguish all animals from all plants. Animals usually lack chlorophyll and the ability to manufacture foods from raw materials available in the soil, water, and atmosphere. Animal cells are usually delimited by a flexible plasma or cell membrane rather than a cell wall composed either of cellulose or chitin, as are the cells of most plants. Animals generally are limited in their growth and most have the ability to move in their environment at some stage in their life history, whereas plants are usually not restricted in their growth and the majority are stationary. The presence or lack of chlorophyll in an organism does not determine its affinity to the plant or animal kingdom. Among the protozoa, the class Phytamastigophora includes animals, such as the euglenids, which have chromatophores containing chlorophyll. These organisms are considered to be animals by zoologists and plants by phycologists. The vestige of a feeding apparatus in these protozoa indicates that they have descended from forms that ingested food particles, that is, animals. Higher parasitic plants and the large plant group Fungi also lack chlorophyll. Another borderline group is the slime molds: the Mycetozoa of zoologists and the Myxomycophyta of the botanists. These organisms exhibit both plant and animal characteristics during their life history. Movement is not a characteristic restricted to the animal kingdom; many of the thallophytes such as Oscillatoria, numerous bacteria, and colonial chlorophytes are motile. The grouping of living organisms into kingdoms is still unsettled. Previously biologists used only two
groups, the Animalia and the Plantae, based on the large, familiar organisms then known to them. The discovery of microorganisms, such as bacteria and viruses, introduced difficulties for this simple scheme. Today biologists recognize up to five kingdoms, but with considerable differences of opinion. Most, however, place the one-celled animals and plants, sometimes along with algae and certain other groups, into the Protista. Other kingdoms are the Monera for the bacteria and blue-green algae, and the Fungi for the slime molds and true fungi. These schemes for recognizing additional kingdoms have the practical advantage of eliminating the difficulties of delimiting and describing the kingdoms of multicellular animals and plants. See ANIMAL KINGDOM; PLANT; PLANT KINGDOM. Walter J. Bock
Animal communication A discipline within the field of animal behavior that focuses upon the reception and use of signals. Animal communication could well include all of animal behavior, since a liberal definition of the term signal could include all stimuli perceived by an animal. However, most research in animal communication deals only with those cases in which a signal, defined as a structured stimulus generated by one member of a species, is subsequently used by and influences the behavior of another member of the same species in a predictable way (intraspecific communication). In this context, communication occurs in virtually all animal species, if only as a means by which a member of one sex finds its partner. The field of animal communication includes an analysis of the physical characteristics of those signals believed to be responsible in any given case of information transfer. A large part of this interest is due to technological improvements in signal detection (for example, the use of tape and video recorders and gas chromatographs), coupled with analysis of the signals obtained with such devices. Communication modes. Information transmission between two individuals can pass in four channels: acoustic, visual, chemical, and electrical. An individual animal may require information from two or more channels simultaneously before responding appropriately to reception of a signal. Furthermore, a stimulus may evoke a response under one circumstance but be ignored in a different context. Acoustic. Sound signals have characteristics that make them particularly suitable for communication, and virtually all animal groups have some forms which communicate by means of sound. Sound can travel relatively long distances in air or water, and obstacles between the source and the recipient interfere little with an animal’s ability to locate the source. Sounds are essentially instantaneous and can be altered in important ways. Both amplitude and frequency modulation can be found in sounds emitted by animals; in some species sound signals have discrete patterns due to frequency and timing of utterances. Since a wide variety of sound signals are
679
680
Animal communication possible, each species can have a unique set of signals in its repertoire. Animals can travel directly toward a sound source, despite variable wind speed and direction, the onset of darkness, or simultaneous signaling by thousands of other animals of various species. Sound signals are produced and received primarily during sexual attraction, including mating and competition. They may also be important in adult–young interactions, in the coordination of movements of a group, in alarm and distress calls, and in intraspecific signaling during foraging behavior. See REPRODUCTIVE BEHAVIOR. The human hearing range (for example, of frequency and intensity) is somewhat limited, but with new technology sounds outside the range of human perception can be tape recorded. Such research reveals that many species of animals, particularly insects, apparently can communicate at frequencies much higher or lower than the human hearing range; a few species, perhaps elephants, communicate at frequencies below the human range. See HEARING (HUMAN); PHONORECEPTION. Visual. Visual signaling between animals can be an obvious component of communication, as is the case for fireflies. Travel time is essentially instantaneous; however, visual signals have very different characteristics from those of sound. Sound requires a medium for transmission (such as air, water, or a substrate), but light travels unimpeded unless blocked by obstacles. Besides the normal range of human vision (visible light), visual signals, that is, the entire electromagnetic spectrum, include additional frequencies in the infrared (heat) and ultraviolet ranges. The quality of light that is often considered is color, but other characteristics are important in visual communication. Alterations of brightness, pattern, and timing also provide versatility in signal composition. Two other features of vision deserve mention. First, an animal equipped with binocular vision can easily locate the precise source of a signal. Second, color patterns in animals (for example, butterflies) can serve as relatively permanent signals. Even with all its favorable characteristics, the visual channel suffers from the important limitation that all visual signals must be line of sight. Information transfer is therefore largely restricted to the daytime (except for animals such as fireflies) and to rather close-range situations because of obstacles normally present in the environment. Intraspecific visual signaling, as with acoustical and chemical signaling, appears to occur primarily during mate attraction. The color dimorphism of birds, the interesting patterns of butterfly wings, the posturing of some fish, and firefly flashing are examples. Among fireflies, for example, flying males flash at regular intervals in specific patterns as they cruise along. Stationary females of the same species, if receptive, flash at a set interval of time after the male, enabling the male to travel directly to them and not to a member of another species. Some parent–young interactions involve visual signaling. A young bird in the nest may open its mouth
when it sees the underside of its parent’s beak. Other examples are the synchronized behavior observed in schooling fish and flocking birds. As with sounds, the human range of perception of visual signals is somewhat limited. Some other species, such as honeybees, see well into the ultraviolet spectrum, and a few forms, such as some beetle species, can perceive objects in the infrared range. Chemical. Chemical signals, like visual and sound signals, can travel long distances, but with an important distinction. Distant transmission of chemical signals requires a movement of air or water. Therefore, an animal cannot perceive an odor from a distance; it can only perceive molecules brought to it by a current of air or water. Animals do not hunt for an odor source by moving other than upwind or upcurrent in water, because chemical signals do not travel in still air or water since diffusion is far too slow. The fact that chemical signals comprise molecules means that they have special attributes important in information transmission, attributes not shared by the other channels. Although acoustical or visual signals are essentially instantaneous, chemical signals have a time lag. A single molecule might stimulate a flying insect to head upwind, but if no other similar molecules are perceived, the insect might well abandon the search or drop downwind until it perceives other molecules of the same type. Chemical signals at times have to be of an appropriate concentration if they are to be effective. A chemical normally considered to be an attractant can serve as a repellent if it is too strong. Chemical signals may persist for a while, and time must pass before the concentration drops below the threshold level for reception by a searching animal. Since molecules of different sizes and shapes have varying degrees of persistence in the environment, the chemical channel is often involved in territorial marking, odor trail formation, and mate attraction to a fixed location. This channel is particularly suitable where acoustical or visual signals might betray the location of a signaler to a potential predator. The array of molecular structure is essentially limitless, permitting a species-specific nature for chemical signals. Unfortunately, that specificity can make interception and analysis of chemical signals a difficult matter for research. Pheromones are chemical signals that are produced by an animal and are exuded through glandular openings (or otherwise) and influence the behavior of other members of the same species. If pheromones are incorporated into a recipient’s body (by ingestion or absorption), they may chemically alter the behavior of such an individual for a considerable period of time. See PHEROMONE. In social insects, long-term regulation of the behavior of individuals within a colony appears to be mainly, if not exclusively, due to chemical inhibition or enhancement of the endocrine system of individuals by pheromones received from other members of the colony. Queen honeybees or queen termites produce one or more chemical substances that inhibit
Animal evolution other females (workers) from becoming functional queens. See SOCIAL INSECTS. All animal species seem to have a somewhat limited range of types of molecules that they can perceive. A substance with a particular odor or taste for humans may not even be perceived by a member of another species. The converse is also true: another species may well be able to smell or taste a compound that humans do not notice. See CHEMICAL ECOLOGY; CHEMORECEPTION. Electrical. Some electric fish and electric eels live in murky water and have electric generating organs that are really modified muscle bundles. Communication by electric signaling is rapid; signals can travel throughout the medium (even murky water), and rather complex signals can be generated, permitting species-specific communication during sexual attraction. However, the electrical mode is apparently restricted to those species that have electric generating organs. See ELECTRIC ORGAN (BIOLOGY). Context. Any of the above types of signals can be effective under the right circumstance but ineffective under other circumstances. For example, when recruited honeybees are searching for a potential hive site found by experienced foragers (scout bees), they are attracted by the chemical components of an exudate from a gland at the upper side and tip end of the abdomen. Similarly, when a honeybee swarm is moving from the location of its parent colony to a new site, an odor from the gland orients bees during transit and stimulates settling once the new site has been reached. However, recruited bees searching for flowers exploited by foragers from their colony apparently ignore those same chemicals. Similarly, male moths of some species are first attracted to the odor of the host plant of their future offspring (plant-feeding caterpillars) and are then attracted to the pheromones exuded by females perched on that host plant. In some cases, if males are not first exposed to the odor of the larval host plant, they may not respond to pheromones exuded by females of their species. See PLANT. Methodology. Animal communication is one of the most difficult areas of study in science for several reasons. First, experiments must be designed and executed in such a manner that extraneous cues (artifacts) are eliminated as potential causes of the observed results. Second, once supportive evidence has been obtained, each hypothesis must be tested. In animal communication studies, adequate tests often rely upon direct evidence—that is, evidence obtained by artificially generating the signal presumed responsible for a given behavioral act, providing that signal to a receptive animal, and actually evoking a specific behavioral act in a predictable manner. A classic example of such a test entails holding a wooden dowel with a red dot on its underside over a gull chick, eliciting the same behavior in the chick as when a parent gull holds its beak over the nest. Whereas this example is quite simple and can be repeated with little equipment, acoustic, visual, chemical, and electric signals can be quite complex, requir-
ing information derived from physics, chemistry, and mathematics and employing electronics and computers. Nonhuman language. The fact that human beings have a language indicates that language is possible in an animal species, and this possibility has been investigated repeatedly in several nonhuman species. Dog, cat, and horse owners, for example, may have a strong bias on this point; those who work with personable animals, such as chimpanzees and dolphins, serve as another example. Experiments conducted by Karl von Frisch in the 1940s yielded results suggestive of a language among honeybees in the form of movements by scout bees through which the location of food sources is communicated to other bees. Although the hypothesis drew wide interest and support because of its emotional appeal, experiments in the 1960s, coupled with a careful analysis of the design and controls of the original series of experiments, indicated that Frisch’s own earlier (1937) odor-search hypothesis more than adequately accounted for the observation. Continued attempts to demonstrate language use in various species (bees, pigeons, chimpanzees, dolphins), and experiments using a computer-simulated bee, have not been able to satisfy critics of the design or interpretation of the experiments. It is possible that human beings may well be the only animal species capable of using symbolic language. See ETHOLOGY; PSYCHOLINGUISTICS. Adrian M. Wenner Bibliography. T. Lewis (ed.), Insect Communication, 1984; T. A. Sebeok (ed.), How Animals Communicate, 1977; A. M. Wenner, Concept-centered vs. Organism-centered Research, 1989; A. M. Wenner, D. Meade, and L. J. Friesen, Recruitment, search behavior and flight ranges of honeybees, Amer. Zool., 31:768–782, 1991; A. M. Wenner and P. H. Wells, Anatomy of a Controversy: The Question of a “Language” among Bees, 1990.
Animal evolution The theory that modern animals are the modified descendants of animals that formerly existed and that these earlier forms descended from still earlier and different organisms. Animals are multicellular organisms that feed by ingestion of other organisms or their products, being unable to derive energy through photosynthesis or chemosynthesis. Animals are currently classed into about 30 to 35 phyla, each of which has evolved a distinctive body plan or architecture; representatives of the major phyla are depicted in Fig. 1. Invertebrates All phyla began as invertebrates, but lineages of the phylum Chordata developed the internal skeletal armature, with spinal column, which was exploited in numerous fish groups and which eventually gave rise to terrestrial vertebrates. The number of phyla is uncertain partly because most of the branching patterns and the ancestral body plans from which
681
coelenterates
chelicerate arthropods
protozoa
flatworms
mollusks
trochophore
minor protostomes
to plants
dipleurula
fishes
echinoderms
primitive chordates
amphibians
reptiles
birds
Fig. 1. Evolution of animal groups, showing hypothetical relationships. ( After G. B. Moment, General Zoology, 2d ed., Houghton Mifflin, 1967)
sponges
ctenophores
nematodes and allies
annelids
mandibulate arthropods
mammals
682 Animal evolution
Animal evolution putative phyla have arisen are not yet known. For example, arthropods (including crustaceans and insects) may have all diversified from a common ancestor that was a primitive arthropod, in which case they may be grouped into a single phylum; or several arthropod groups may have evolved independently from nonarthropod ancestors, in which case each such group must be considered a separate phylum. So far as known, all animal phyla began in the sea. See ANIMAL; ANIMAL KINGDOM. Several lines of evidence bear significantly upon the ancestry of these animals and of their major subdivisions. Classically, animals were grouped according to the similarities of their adult body plans, and while the body plans have provided a basis for the recognition of most of the phyla and other distinctive animal groups, they have not provided definitive evidence of the interrelations of the phyla. A second approach is to compare patterns of early development, which should be more conservative than adult features; such comparisons have permitted the grouping of some phyla into likely alliances. Hypotheses as to the body plans of phylum ancestors have been erected from these morphologic data. It has not proven possible to corroborate any hypotheses, however, for the earliest members of the phyla (and of invertebrate classes) appear abruptly in the fossil record, and their ancestors cannot be traced. Furthermore, numbers of extinct phyla or other major animal groups also appear in the fossil record, adding branches to the tree of life which must be reckoned with, but with no definitive indication of either their origins or branching patterns. Finally, genetic changes that have accumulated after branching events separated the major groups should have left a pattern of variation in the structure of genes that may be reconstructed to yield a phylogeny. If genetic changes have occurred in a reasonably regular manner, then the amount of divergence in the structure of genes and gene products, such as protein and ribonucleic acid (RNA) molecules, should be proportional to the time since branching occurred. Some genes change so fast that their pattern of divergence is useful only for recent branching events, while others are so conservative that they are used to study divergences that occurred billions of years ago. These methods are promising when used in conjunction with other evidence. Major animal groups. The simplest living phyla may have no direct relationships to other animals (Fig. 2), but indicate the versatility of multicellular design. Porifera (sponges) contain several differentiated cell types that form body walls, but their architecture is limited to folding these walls into chambers that are sometimes highly convoluted. Coelenterata (=Cnidaria; sea anemones, corals, and jellyfish) have two cellular body layers, separated by a gelatinous layer; the inner tissue forms a digestive cavity. See CNIDARIA; PORIFERA. Flatworms (Platyhelminthes) and their allies have solid bodies, with both external (ectoderm) and internal (endoderm) tissue layers separated by a third cellular layer (mesoderm); there is a gut lined
Chordata
Annelida
Urochordata Sipuncula
Bryozoa Phoronida
Hemichordata
Pogonophora
Brachiopoda
arthropod phyla
Echinodermata
Mollusca Nemertina Platyhelminthes Porifera
Cnidaria Protozoa (Flagellata)
Fig. 2. Phylogenetic tree depicting the possible relations among the phyla according to evidence from all sources. The branching pattern is constrained by models of development and of body-plan evolution and by molecular data. Branch length is not to scale.
with endoderm. All phylogenetic evidence suggests that a flatwormlike form gave rise to the remaining phyla, which possess body cavities (other than a gut) in which reproductive and other organs can be sequestered. In about 15 of these phyla this cavity (the coelom) is lined by mesoderm and usually functions as a hydrostatic skeleton. See PLATYHELMINTHES. The coelomates constitute all of the familiar, largerbodied animals and are commonly grouped into two main divisions (Fig. 2). In one group, the body was originally segmented, as in annelid worms or arthropods, though the segmentation may be reduced or lost in some branches. In the other group, the coelom is divided into two or three regions; these include the Enchinodermata (starfish and sea urchins), Tunicata (sea squirts), and Chordata (including humans). Additionally, there are coelomate phyla that cannot yet be assigned to either of these groups; some, such as the Brachiopoda, Phoronida, and Bryozoa, display an enigmatic mixture of traits that otherwise characterize one or the other of the main groups. In about 10 phyla (not shown in Fig. 2) the body cavity is not a true coelom but lies between endoderm and mesoderm and is termed a pseudocoel. These pseudocoelomates are soft-bodied, small to minute, and commonly parasitic; some of them have never been found as fossils. See ANNELIDA; BRACHIOPODA; BRYOZOA; CHORDATA; ECHINODERMATA; PHORONIDA; TUNICATA. Origins of animals. Some features of the cells of primitive animals resemble those of the single-celled Protozoa, especially the flagellates, which have long been believed to be animal ancestors. Molecular phylogenies have supported this idea and also suggest that the phylum Coelenterata arose separately from all other phyla that have been studied by this technique. Thus animals may have evolved at least twice from organisms that are not themselves animals, and represent a grade of evolution and not a single branch (clade) of the tree of life. Sponges have also been suspected of an independent origin, and it is possible
683
684
Animal evolution that some of the extinct fossil phyla arose independently or branched from sponges or cnidarians. See PROTOZOA. The earliest undoubted animal fossils (the Ediacaran fauna) are soft-bodied, and first appear in marine sediments nearly 650 million years (m.y.) old. This fauna lasted about 50 m.y. and consisted chiefly of cnidarians or cnidarian-grade forms, though it contains a few enigmatic fossils that may represent groups that gave rise to more advanced phyla. Then, nearly 570 m.y. ago, just before and during earliest Cambrian time, a diversification of body architecture began that produced most of the living phyla as well as many extinct groups. The body plans of some of these groups involved mineralized skeletons which, as these are more easily preserved than soft tissues, created for the first time an extensive fossil record. The earliest mineralized skeletons are chiefly minute, the small shelly fauna of earliest Cambrian age; and while some of these fossils are primitive members of living phyla, many—perhaps 10 to 20 kinds of them—are distinctive and quite possibly represent phyla or at least major branches of phyla that soon became extinct, geologically speaking. The soft-bodied groups were also markedly diversified, though their record is so spotty that their history cannot be traced in detail. A single, exceptionally
65 235
a
bit
recent C M
lo Tri
ora ora a a ph ph o ata o a a t t hor ea hin c on op rinoid stom hozo olit pla ticula litha d n t o o o p o p n r o r n c g a y o a y r e o o o c H E In P G M C M S
ta ae
h
lyc
Po
preserved soft-bodied fauna from the Burgess Shale of British Columbia that is about 530 m.y. old contains not only living soft-bodied worm phyla, but extinct groups (perhaps a dozen) that cannot be placed in living phyla and do not seem to be ancestral to them. Among these is Wiwaxia, a novel form which bears spinelike elements similar to some of the small shelly fossils of Early Cambrian age, thus providing evidence that the small shelly forms did indeed include novel types. Clearly, an early radiation of animals produced a vast array of novel body plans in an evolutionary episode that has never been surpassed. See BURGESS SHALE; CAMBRIAN; EDIACARAN BIOTA; FOSSIL. History of major groups. Following the early phase of rampant diversification and of some concurrent extinction of phyla and their major branches, the subsequent history of the durably skeletonized groups can be followed in a general way in the marine fossil record (Fig. 3). The composition of the fauna changed continually, but three major associations can be seen: one dominated by the arthropodlike trilobites during the early Paleozoic (Fig. 3a), one dominated by articulate brachiopods and crinoids (Echinodermata) in the remaining Paleozoic (Fig. 3b), and one dominated by gastropod (snail) and bivalve (clam) mollusks during the
P 570 680
pC
million years ago
(a)
C M
Art
a ide
no
ic
recent 65 235
ta ula
Cri
c
tra
Os
da
a od
o lop
pha
a ozo
th
An
Ce
ra a ia ho ata dea ng na op omat ea m t a i e n ithi ospo d t o l a i o l r ide s o o d o e t t r o ll n ro las p le sto Me B Ste Ste Con Sc Cy Gra
P 570 680
pC C p (b)
da
o op str
recent 65 235
C M
Ga
ia
alv
Biv
aca
es
hy cht
tei
Os
r ost
lac
Ma
es a ta ia ng ichthy inellid idea a ma o e p a oa s re ro ct dr ol mo hon exa telle alca ydroz mn H C C S De Gy H
a
hin
e oid
Ec
P 570
pC
680 (c)
Key:
C = Cenozoic
M = Mesozoic
P = Paleozoic
pC = Precambrian
Fig. 3. Marine fossil record of major groups during the last 680 million years. The width of the spindles is proportional to the number of families present. The shifting of dominance from (a) trilobites and other forms during the early Paleozoic to (b) articulate brachiopods, crinoids, and others in the late Paleozoic and then to (c) gastropod and bivalve mollusks can be seen. (After J. J. Sepkoski, Jr., in J. W. Valentine, ed., Phanerozoic Diversity Patterns: Profiles in Macroevolution, Princeton University Press, 1985)
Animal evolution Mesozoic and Cenozoic (Fig. 3c). The mass extinction at the close of the Paleozoic that caused the contractions in so many groups (Fig. 3) may have extirpated over 90% of marine species and led to a reorganization of marine community structure and composition into a modern mode. Resistance to this and other extinctions seems to have been a major factor in the rise of successive groups to dominance. On land, the fossil record is too spotty to trace in detail the invasions of the terrestrial environment, which must have been under way by the mid-Paleozoic. Annelids, arthropods, and mollusks are the more important invertebrate groups that made the transition to land. The outstanding feature of terrestrial fauna is the importance of the insects, which appeared in the late Paleozoic and later radiated to produce the several million living species, surpassing all other life forms combined in this respect. See ANNELIDA; ARTHROPODA; CENOZOIC; INSECTA; MESOZOIC; MOLJames W. Valentine LUSCA; PALEOZOIC.
itive nonvertebrates, the protochordates: lancelets, tunicates, acorn worms, pterobranchs, and possibly the extinct graptolites and conodonts. The interrelationships of these forms are not well understood. With the exception of the colonial graptolites, they are soft-bodied and have only a very limited fossil record. They suggest possible links to the Echinodermata in developmental, biochemical, and morphological features. In addition, some early Paleozoic fossils, the carpoids, have been classified alternatively as chordates and as echinoderms by various students, again suggesting a link. In spite of these various leads, the origin of the chordates remains basically unclear. See VERTEBRATA. Chordates are characterized by a hollow, dorsal, axial nerve chord, a ventral heart, a system of slits in the pharynx that serves variously the functions of feeding and respiration, a postanal swimming tail, and a notochord that is an elongate supporting structure lying immediately below the nerve chord. The protochordates were segmented, although sessile forms such as the tunicates show this only in the swimming, larval phase. Even the free-swimming forms were not truly active, and they used the gill apparatus primarily for feeding on small particles
Chordate Origins
300 Carboniferous
Ordovician
Synapsida
Diplorhina
Silurian
450
?
500 Cambrian 550 Fig. 4. Vertebrate evolution in relation to geologic time scale. Classes are shown in capital letters; subclasses are shown in lowercase letters, except that “Hagfishes” and “Lampreys” are common names for orders. (After M. Hildebrand, Analysis of Vertebrate Structure, Wiley, 1982)
Theria
Prototheria
MAMMALIA
Neornithes
Archaeornithes
Archosauria
Lepidosauria
Euryapsida
Anapsida
R E P T I L I A
Devonian Monorhina
400
Paleozoic
350
Labyrinthodontia
AGNATHA
Lepospondyli
Permian
PLACODERMI
250
AMPHIBIA
Sarcopterygii
Triassic
ACANTHODII
CHONDRICHTHYES 200
Lissamphibia
OSTEICHTHYES
Actinopterygii
Jurassic
Elasmobranchii
Mesozoic
AVES
Tertiary
Cretaceous
100
150
Quaternary
Holocephali
50
Evolution of the Vertebrate Classes and Subclasses
Periods
Lampreys
Eras
Hagfishes
Millions of years
Cenozoic
The phylum Chordata consists largely of animals with a backbone, the Vertebrata, including humans (Fig. 4). The group, however, includes some prim-
685
686
Animal evolution trapped on mucous streams generated in the pharynx. There are two hypotheses concerning the nature of the protochordates. One holds that all protochordates were originally sessile, with, however, active larvae, and that the vertebrates arose from the larval phase through increased development of the segmented muscular trunk and associated nervous system. The alternative hypothesis proposes that the sessile condition seen in many protochordates is secondary. The first vertebrates were fishlike animals in which the pharyngeal slits formed a series of pouches that functioned as respiratory gills. An anterior specialized mouth permitted ingestion of food items large in comparison with those of the filter-feeding protochordates. Vertebrates are first known from bone fragments found in rocks of Cambrian age, but more complete remains have come from the Middle Ordovician. Innovations, related to greater musculoskeletal activity, included the origin of a supporting skeleton of cartilage and bone, a larger brain, and three pairs of cranial sense organs (nose, eyes, and ears). At first the osseous skeleton served as protective scales in the skin, as a supplement to the notochord, and as a casing around the brain. In later vertebrates the adult notochord is largely or wholly replaced by bone, which encloses the nerve chord to form a true backbone. All vertebrates have a heart which pumps blood through capillaries, where exchanges of gases with the external media take place. The blood contains hemoglobin in special cells which carry oxygen and carbon dioxide. During the exchanges the oxygen content of the cells is increased and that of carbon dioxide decreased. In most fishes the blood passes from the heart to the gills and thence to the brain and other parts of the body. In most tetrapods, and in some fishes, blood passes to the lungs, is returned to the heart after oxygenation, and is then pumped to the various parts of the body. Fish evolution. The jawless fish, known as Agnatha, had a sucking-rasping mouth apparatus rather than true jaws. They enjoyed great success from the Late Cambrian until the end of the Devonian. Most were heavily armored, although a few naked forms are known. They were weak swimmers and lived mostly on the bottom. The modern parasitic lampreys and deep-sea scavenging hagfish are the only surviving descendants of these early fish radiations. See DEVONIAN; JAWLESS VERTEBRATES. In the Middle to Late Silurian arose a new type of vertebrate, the Gnathostomata, characterized by true jaws and teeth. They constitute the great majority of fishes and all tetrapod vertebrates. The jaws are modified elements of the front parts of the gill apparatus, and the teeth are modified bony scales from the skin of the mouth. With the development of jaws, a whole new set of ecological opportunities was open to the vertebrates. Along with this, new swimming patterns appeared, made possible by the origin of paired fins, forerunners of which occur in some agnathans. See GNATHOSTOMATA; SILURIAN.
Four groups of fishes quickly diversified (Fig. 2). Of these, the Placodermi and Acanthodii are extinct. The Placodermi were heavily armored fishes, the dominant marine carnivores of the Silurian and Devonian, rivaled only by the large eurypterid arthropods. The Acanthodii were filter-feeders mostly of small size. They are possibly related to the dominant groups of modern fishes, the largely cartilaginous Chondrichthyes (including sharks, rays, and chimaeras) and the Osteichthyes (the higher bony fishes). These also arose in the Late Silurian but diversified later. See ACANTHODII; CHONDRICHTHYES; OSTEICHTHYES; PLACODERMI. Conquest of land. The first land vertebrates, the Amphibia, appeared in the Late Devonian and were derived from an early group of osteichthyans called lobe-finned fishes, of which two kinds survive today, the Dipnoi or lungfishes, and the crossopterygian coelacanth Latimeria. They were lung-breathing fishes that lived in shallow marine waters and in swamps and marshes. The first amphibians fed and reproduced in or near the water. True land vertebrates, Reptilia, with a modified (amniote) egg that could survive on land, probably arose in the Mississippian. See AMNIOTA; AMPHIBIA; DIPNOI; MISSISSIPPIAN. Reptile radiations. By the Middle Pennsylvanian a massive radiation of reptiles was in process. In it can be traced the line leading to mammals as well as the lineages of the great reptiles that dominated the Mesozoic Era. The most prominent reptiles belong in the Diapsida: dinosaurs, lizards and snakes, and pterosaurs (flying reptiles). The birds, Aves, which diverged from the dinosaur radiation in the Late Triassic or Early Jurassic, are considered to be feathered dinosaurs, and thus members of the Diapsida, whereas older authorities prefer to treat them as a separate class. In addition, there were several Mesozoic radiations of marine reptiles such as ichthyosaurs and plesiosaurs. Turtles (Chelonia) first appeared in the Triassic and have been highly successful ever since. See AVES; DINOSAURIA; JURASSIC; PENNSYLVANIAN; REPTILIA. Mammalian origins. The line leading to mammals can be traced to primitive Pennsylvanian reptiles, Synapsida, which diversified and spread worldwide during the Permian and Triassic. The first true mammals, based on characteristics of jaw, tooth, and ear structure, arose in the Late Triassic. Derived mammals, marsupials (Metatheria) and placentals (Eutheria), are known from the Late Cretaceous, but mammalian radiations began only in the early Cenozoic. By the end of the Eocene, all the major lines of modern mammals had become established. Mammals are easily separated into distinct groups (orders), but their relationships are not easy to discover from fossil records because of the explosion of mammalian evolution in the early Cenozoic. Molecular analyses (blood proteins, deoxyribonucleic acid, ribonucleic acid) of living mammals show that the most primitive group of placentals is the edentates (sloths, armadillos, and anteaters). An early large radiation included the rodents, primates (including monkeys, apes, and
Animal feeds humans), and bats, possibly all closely related to the insectivores and carnivores. The newest radiations of mammals are of elephants and sea cows, while the whales are related to the artiodactyls (cattle, camels). Because of biogeographic isolation, marsupials came to flourish in Australia and South America, whereas the placentals diversified widely in Eurasia and North America. See CENOZOIC; CRETACEOUS; EOCENE; MAMMALIA; PERMIAN; SYNAPSIDA; TRIASSIC. Primates and humans. Among the early and primitive lines of placental mammals were the Paleocene and Eocene plesiadapoids from which the primates arose. They were mostly nocturnal, insect- and fruiteating animals with forwardly directed eyes and a locomotor system well developed for an arboreal life. The living primates consist of prosimians; tarsiers; lemurs; monkeys and anthropoids; apes; and humans and near relatives, the hominids. During the Oligocene, monkeys diverged into Old and New World lines. From the former, the apes arose in the late Oligocene. Known as the Pongidae, they are represented today by gibbons, orangutans, chimpanzees, and gorillas. See APES; OLIGOCENE. The Hominidae, which include the human subspecies Homo sapiens sapiens, diverged from the apes about 5 million years ago. The closest living relatives among the apes are chimpanzees and gorillas. The earliest known hominids consist of fossils of the genus Australopithecus, excavated from rocks formed about 3.75 million years ago in eastern Africa. Several distinct lines of Australopithecus, under various generic names, coexisted for a long period of time; and also, during the later part of their existence, with H. habilis, which lived from about 2 million to 1.75 million years ago. This early species of Homo was notable for the relatively large brain, and probably included the first tool user. It was superseded by H. erectus, a larger animal with a greatly increased relative brain size. About 1 million years ago, H. erectus spread from Africa to eastern and southern Asia. It too has been given several generic names, of which Pithecanthropus and Sinanthropus from eastern Asia are the most familiar. See AUSTRALOPITHECINE; FOSSIL HUMANS. Skeletal remains anatomically similar to those of H. sapiens have been found in northern Africa in beds formed about 300,000 years ago. Truly modern forms appeared some 100,000 years ago. They penetrated Europe some 30,000 to 40,000 years ago and spread throughout the Earth, occupying the Old World first, penetrating to Australia, and about 12,000 to 15,000 years ago entering North America and spreading rapidly to South America. See PALEONTOLOGY. Keith S. Thomson Bibliography. R. L. Carroll, Vertebrate Paleontology and Evolution, 1988; K. G. Field et al., Molecular phylogeny of the animal kingdom, Science, 239: 748– 752, 1988; M. F. Glaessner, The Dawn of Animal Life, 1984; A. S. Romer, Vertebrate Paleontology, 3d ed., 1966; J. W. Valentine (ed.), Phanerozoic Diversity Patterns: Profiles in Macroevolution, 1985; H. B. Whittington, The Burgess Shale, 1985; J. Z. Young, Life of the Vertebrates, 3d ed., 1991.
Animal feeds Substances suitable for the nutrition of animals. This article discusses the main category—livestock and poultry feed— then briefly describes the secondary area—pet foods. Livestock and Poultry Feed In formulating livestock and poultry diets, animal nutritionists select from a variety of bulk feed ingredients including forages, feed grains, protein concentrates, by-product feeds, and vitamin-mineral supplements. These are formulated into complete rations in modern feed manufacturing plants (see illus.) to provide a balance of energy, protein, vitamins, and minerals to meet the nutrient requirements of various classes of livestock and poultry. Drugs and other additives may also be needed to meet specific needs. Today most feed manufacturers use high-speed electronic computers and various linear programming techniques to formulate diets. The overall goal is to prepare the lowest-cost ration that will provide economical livestock and poultry products for the consumer while maximizing returns to the producer. Forages. In a broad sense, forages consist of grasses and legumes fed to ruminants (cattle, sheep, goats) and horses. Because of symbiotic microbes (bacteria and protozoa) in the digestive tract of these animals, they are uniquely equipped to utilize nutrients in forages; thus they do not compete with humans for the same food source. See FORAGE CROPS. Grazing. Pasture and rangeland areas are covered with forage that is harvested by the grazing animal. Pasture provides about one-third of the nutrients consumed by dairy cattle and about half of those used by beef cattle and sheep. Cattle and sheep on rangeland obtain more than half their feed from forage. Grazing quality forage has the advantages of providing more digestible protein per acre than most other feeds, eliminating harvesting and feeding costs, and reducing soil erosion that results from the action of water and wind. While a portion of the forage is returned to the land as manure, appreciable forage may be lost due to trampling and selection by the grazing animals. To preclude such loss, forages are often mechanically harvested and fed as hay, silage, or dehydrated material. Although this practice minimizes loss and permits forage to be harvested at the optimum stage of maturity and highest quality, additional labor and energy are required to operate the harvesting, processing, and feeding equipment. See GRASS CROPS. Hay and silage. In the United States, cool-season grasses, such as tall fescue and orchard grass, grow most rapidly during the spring and fall when rainfall is adequate and temperatures are moderate. Conversely, warm-season species, such as bahia grass, bermuda grass, and digit grass, grow best during the summer. Furthermore, because of variable climatic conditions, forage species do not grow at a uniform rate throughout the growing season. Because of this uneven distribution in forage production, most farms produce more forage than can be consumed by the
687
688
Animal feeds two-stage pneu-vac
primary collector secondary collector air lock
two-way valves
turnhead
3
4
5
6
7
bulk ingredient bins
scale truck dump
surge hopper
bag dump
8 9 10 11 12
to truck bin
belt conveyors
cooler feeder and scalper
bin bin
bin
surge hopper
6 7
collector crumbles rolls
patch mixer
1 2 3 4 5
12-position switch station
collector 1 2
bulk truck bins
high mol. pellet mill duster
bagging scale
bin
pellet mill
sewing belt scalper
scale
truck scale
to cake press cooler
Key:
fines
screw conveyor vertical screw
blower air lock
elevator
Flow diagram of formula feed plant. (Sprout Waldron and Co., Inc.)
livestock. This surplus forage is conserved as hay or silage and either fed or sold at later periods. Hay is dried forage which has been cut and suncured in the field to about 10–20% moisture and then baled or otherwise packaged. It is commonly fed to ruminants and horses when grazing is inadequate or when the animals are confined. In the United States, alfalfa is the major hay crop, but other legumes, as well as annual and perennial grasses, are also extensively grown. The type of hay produced in a region depends primarily on those species which are best adapted to that area. For example, in the Northeast, timothy and alfalfa are important hay crops, while in the humid Southeast, warm-season perennial grasses such as Coastal bermuda grass and Pensacola bahia grass predominate. Silage is fermented forage which is produced by the action of acetic and lactic acid–producing anaerobic bacteria in closed structures called silos. Although almost any forage crop can be preserved as silage, corn ranks first in the United States because good-quality, chopped, whole corn plant ensiles well and produces more digestible nutrients per acre than any other forage. Other silage consists of sorghum, grass, legumes, small grains, and miscellaneous crops. See CLOVER; LEGUME FORAGES. Dehydration. Preserving forage by dehydration is relatively new compared to hay-making or ensiling. Although this industry is based largely on alfalfa, other crops, including Coastal bermuda grass, Johnson grass, and whole corn plant are also dehydrated. In this process, freshly harvested or field-wilted
forage is chopped and fed directly into the combustion gases (1600–2000◦F or 870–1090◦C) inside rotating horizontal drum dryers. These may be single- or triple-pass and range 8–12 ft (2.4–3.7 m) in diameter and 24–65 ft (7–20 m) in length. The evaporative capacity ranges 3–30 tons (2.7–27 metric tons) of water per hour. During drying, moisture diffuses to the surface of the forage particles as fast as it is evaporated. This rapid evaporation of surface moisture keeps the plant material cool enough to prevent burning. After remaining in the drum for about 2–10 min, the dehydrated forage is separated from the moisture-laden exit gases (250–350◦F or 120– 180◦C) by means of a cyclone separator or other type of separator. Next, the forage is ground in a hammer-mill and pelleted for ease of storage, transportation, and feeding. To control dust, about 1.0% animal fat or vegetable oil may be added prior to grinding. Mechanically dehydrated forages such as alfalfa and Coastal bermuda grass usually do not compete with hay or silage as roughages for ruminants. These products are high in protein, minerals, carotene (provitamin A), tocopherol (vitamin E), vitamin K, xanthophylls (substances that provide yellow color to poultry skin and egg yolks), and unidentified growth and reproduction factors. They are used primarily as supplements in swine and poultry rations and pet foods. The importance of carotene, xanthophylls, and vitamin E is such that producers of pelleted, dehydrated alfalfa or Coastal bermuda grass usually store the product under inert gas or add an
Animal feeds antioxidant (ethoxyquin) to minimize oxidation and ensure a high-quality product. Feed grains. Feed grains are any of several grains commonly used as livestock and poultry feeds; they include corn, barley, oats, and grain sorghum (milo). Although wheat is second only to corn as a cereal grain in the United States, wheat is not usually fed to livestock. Most of the wheat crop is used for the manufacture of flour and other human foods. However, when properly used, wheat is a satisfactory feed for all classes of livestock. Accordingly, when the price of feed grains is high and surplus wheat is available at low price, considerable amounts may be fed to animals. See BARLEY; CORN; OATS; SORGHUM; WHEAT. The purpose of processing feed grains is to improve their palatability, digestibility, and overall feed efficiency. For many years, stockmen sought to do this by grinding, crushing, rolling, and soaking feed grains, but since the late 1950s several new processing techniques have been used. One such technique is to ensile high-moisture (32%) ground ear corn. This has been shown to result in 10–15% greater utilization by cattle compared to regular ground ear corn. Additional techniques for improving grains by processing include reconstitution (moistening), flaking, roasting, micronizing, popping, and exploding. Reconstitution. The process in which water is added to grain to bring its moisture content up to 20–30% is known as reconstitution. Usually grain so treated is held in limited-oxygen storage for a minimum of 3 weeks. Generally, cattle fed reconstituted corn or milo gained 8–15% more weight than those fed regular grain; and the feed efficiency of the reconstituted grain is increased by 7–11%. Flaking. In a modification of steam rolling, the grain is cooked or steamed at approximately 200◦F (93◦C) at atmospheric pressure for 15 to 30 min, or at higher temperature and pressure, approximately 300◦F (150◦C) and 50 lb/in.2 (345 kilopascals) for 1 to 2 min, and then is rolled into flakes 1/√32 in. (0.8 mm) thick. Because this process increases the moisture content to about 15–20%, the rolled flakes are immediately dried to 15% moisture. Flaking was the first modern process whereby the rate of gain and feed efficiency of milo were markedly increased. For beef cattle, gain and feed efficiency increased up to about 11% and 15%, respectively. Roasting. In roasting corn, the grain is heated to 300◦F (150◦C). This results in expansion of the kernel, as evidenced by the fact that the roasted material weighs 39 lb/ft3 (625 kg/m3) compared to 45 lb/ft3 (720 kg/m3) for the raw corn material. Moisture content also is decreased about 2–4%. For fattening cattle, roasted corn improves weight gain by about 11% and has a feed efficiency about 14% greater than that of unroasted corn. Micronizing. Another method used in the feed industry involves heating grain to 300◦F (150◦C) with microwaves. This micronized grain contains about 7% moisture. It is then rolled to form a uniform, freeflowing product. Prior to feeding, water is usually added to adjust the moisture content to 10%. Cattle fed micronized grain sorghum (milo) compare favor-
ably with those fed steam-flaked grain sorghum in rate of gain and feed efficiency. However, the cost of processing favors the micronizing technique over the steam-flaking method. Popping. This method involves rapidly heating grain containing about 15–20% moisture to 300–310◦F (150–155◦C). At this temperature, the internal moisture of the grain is volatilized and explodes the kernel. The exploded product is then dry-rolled to form a uniform product and reduce storage space. All grains can be processed by this method; however, it appears to be particularly effective in processing milo. Feeding trials conducted in Texas revealed that cattle consumed less of the popped milo than the regular grain. Although this resulted in lower average daily gains, feed efficiency was increased 17%. Explosion puffing. In another method, explosion puffing, live steam is injected into high-tensile-strength steel bottles which hold about 200 lb (90 kg) of grain. About 20 s after the pressure reaches 250 lb/in.2 (1.7 megapascals) a valve is opened. Rapid release of the pressurized moisture within the grain explosively puffs the kernels. In practice, tandem vessels are alternately filled and pressurized to achieve a continuous operation. The resulting product resembles puffed breakfast cereal. Cattle fed puffed milo have exhibited feed intake, gain, and feed efficiency similar to that of cattle fed flaked grain. Protein supplements. Protein supplements are feeds which contain more than 20% protein or protein equivalent. They are primarily added to rations containing feed grains or forages which are low in protein. Addition of protein supplements depends upon the animal’s requirement for protein and the amount and quality of grain and forage in the ration. Plant sources. High-protein meals can be made from soybean, cottonseed, sunflower seed, safflower seed, rapeseed, peanut, linseed, and coconut (copra). These meals are by-products of the oilseed-crushing industry; they vary in feeding value (that is, protein quantity and quality) depending upon the amount of hull or seedcoat remaining in the meal, conditions of time and temperature during extraction, and constituent amino acids of the protein. For example, although soybean meal and sunflower meal contain about the same amount of protein (45.8% versus 46.8%, respectively), they differ in amino acid content. Compared to sunflower meal, soybean meal is high in lysine (6.6 versus 3.8 g/16 g nitrogen) and low in methionine (1.1 versus 2.1 g/16 g nitrogen). In the United States, soybean meal is the most widely used protein supplement. Quality forage and other leafy plants contain considerable protein (12–30%). The amino acid content of this protein compares favorably to that of soybean meal, fishmeal, and other high-quality proteins. Considerable research has been conducted in the United States, England, and Pakistan to investigate the preparation of leaf protein concentrates from forages. In the United States, the Pro-Xan process was developed for concentrating protein from alfalfa. The final product constitutes 4.5% of the starting material and
689
690
Animal feeds contains 40–50% protein and 600–900 mg/kg xanthophyll. Feeding trials with chicks have indicated that the product is a good source of protein and xanthophyll. If the cost of processing can be lowered, leaf protein concentrate may contribute significantly as an alternate protein source in the production of animal protein—meat, milk, and eggs—for humans. Animal sources. Protein supplements of animal origin primarily consist of inedible products derived from meat packing or rendering plants, marine sources, poultry processing plants, and dairies. In general, such products are high in protein and of excellent quality, being well balanced in amino acids and providing vitamins and minerals. Feather meal, a by-product of the poultry processing industry, has also been used as a protein source. Although this product contains about 85% protein, it is poorly digested and must be hydrolyzed for good utilization. Animal-derived protein supplements are subject to autoxidation and rancidity because of their high fat content, and are often a source of bacterial (particularly Salmonella) contamination of poultry and meat products used for human food. It is therefore important that proper sanitation be practiced during processing and handling of these protein supplements to prevent bacterial contamination of human food. Single-cell sources. There has been considerable interest in the production of protein from single-cell organisms such as yeasts, bacteria, and algae. A wide variety of carbon sources can be used as substrate for the growth of these organisms, including petroleum, methane, alcohol, starch, molasses, cellulose, and paper pulp waste liquor. Potential yields are quite high. One thousand pounds (450 kg) of single-cell organisms might produce 50 tons (45 metric tons) of protein per day. Although production of single-cell protein as a feedstuff for livestock appears promising, problems relating to palatability, digestibility, nucleic acid content, toxins, protein quality, and economics must be solved before this becomes a reality. Microbial synthesis. Cattle and other ruminants are able to derive a portion of their dietary protein from nonprotein nitrogen via microbial protein synthesis in the rumen. Amino acids, amides, ammonium salts, and other nonprotein nitrogen sources provide nitrogen from which symbiotic rumen bacteria synthesize microbial protein. This microbial protein is then digested by the ruminant animal in the abomasum and gastrointestinal tract. Providing that adequate energy and other nutrients are available, about one-fourth to one-third of the protein requirements of ruminants can be met by the feeding of nonprotein nitrogen. If this amount is exceeded and the animals are not managed properly, toxicity and death may result. Commercially available products include urea and the ammoniated products of molasses, citrus pulp, beet pulp, cottonseed meal, and rice hulls. Both solid (range cubes and pellets) and liquid (lick tanks) products containing urea, molasses, trace minerals, vitamin A, and so forth, are available. Slow-release urea products also are marketed. With
concern about world food supply and competition between humans and animals for available protein, the feeding of nonprotein nitrogen to ruminants is increasingly important. See PROTEIN. Vitamins. Vitamins are complex organic compounds which are required in exceedingly small amounts by one or more animal species. Most function as cofactors in enzyme systems involving the regulation of metabolism and transformation of energy. Certain vitamins (such as vitamin K) are apparently needed by only a few species; others (such as vitamin A and thiamine) are required by all animals. Some vitamins (such as vitamin A) must be present in the diet; others are synthesized in the tissues of the animals (such as vitamin C) or by bacteria in the digestive tract (such as B-complex vitamins). The absence in the diet of one or more of the required vitamins may prevent the animal from growing or reproducing, or may cause deficiency diseases. If the deficiency is severe, the animals may die. Farm animals fed diets containing high-quality green forage usually receive enough of most required vitamins. Because monogastric animals (poultry, swine) thrive when small amounts of leaf meal are added to their diets, alfalfa and other forages are usually described as containing unidentified growth factors. High temperatures and oxidation destroy certain vitamins; therefore, care must be taken in the processing and storage of feedstuffs to protect the vitamins they contain. Certain feedstuffs also may contain antivitamins, which prevent the proper assimilation and functioning of vitamins. To ensure that rations contain sufficient vitamins to meet dietary requirements, manufacturers often fortify feed formulations with one or more vitamins. See VITAMIN. Minerals. Animals, like plants, require minerals to grow, and if deficiencies occur, livestock become unthrifty, lose weight, or exhibit other deficiency symptoms. Minerals serve many vital functions. The skeletons of vertebrates are composed chiefly of calcium and phosphorus. Not only must livestock receive adequate calcium and phosphorus, but they must also receive these minerals in the proper proportion, because a great excess of one or the other may be detrimental. Also, if the proper calcium-phosphorus ratio is provided, less vitamin D is needed. Minerals also are necessary constituents of soft tissues and body fluids, and they function in maintaining the proper acid-base balance of body tissues. In addition to macroelements (for instance, calcium, phosphorus, sodium, chlorine, and iron), trace elements (zinc, cobalt, manganese, and selenium, among others) are required for proper nutrition. Trace elements primarily function as cofactors in metabolic reactions. Feedstuffs produced in particular areas may be deficient in or have an excess of certain minerals. For example, in the Great Lakes and Northwest regions of the United States, soils are deficient in iodine. Also, soils containing more than 0.5 ppm (part per million) of selenium are potentially dangerous for production of livestock feed. The form and availability of certain minerals also are an important consideration in formulating livestock diets. Phytin phosphorus,
Animal feeds for example, is not a good source of phosphorus for poultry or swine; however, it appears to be satisfactory for ruminants. To ensure an adequate supply of minerals, feed manufacturers often add mineral supplements. Farm animals may also be permitted free access to a mineral mixture. If this is allowed, care should be taken, particularly with trace minerals, to see that the animals do not consume too much, or toxicity may result. Like vitamins, use of mineral supplements adds to production costs. Therefore, the need for such supplements should be given careful consideration. By-products. This group of feeds includes a number of by-products formed in the processing of livestock, plants, and other materials. To be acceptable as livestock feeds, such products should be palatable, provide digestible nutrients or roughage, and present no health hazard. Use of by-products as feed not only lowers the cost of animal production but helps to reduce the competition between livestock and humans for available grain and protein supplies. Another advantage is that industry derives income from products that otherwise might be wasted and disposed of at a cost. By-product feeds include inedible products and residues from animal, poultry, and marine processing plants, fruit and vegetable processing plants (citrus pulp, tomato pomace), the brewing and distilling industry (brewers’ and distillers’ grains), dairy processing plants (whey, cheesemeal), wood and paper mills (wood molasses, treated wood scraps), as well as miscellaneous fibrous materials such as cottonseed gin trash, cottonseed, soybean, and peanut hulls. Research has been conducted on the feeding of animal wastes (manure and litter) to livestock. There is also interest in the treatment of crop residues (cornstalks, straw, and so forth) with alkali (sodium hydroxide) to improve their digestibility for ruminants. Although certain by-products (such as peanut hulls) are available only in a limited area, others (such as brewers’ grains) are more widely available. Feed additives and implants. These may be defined as nonnutritive substances which enhance the utilization of a feed or productive performance of the animal. It is estimated that in the United States 75% of the finishing cattle, finishing lambs, and growingfinishing pigs receive feed additives or implants. Use of these products and guidelines for administration and withdrawal of approved products are regulated and prescribed by the U.S. Food and Drug Administration (FDA). Also, feed-control officials in individual states may regulate use of these substances for livestock marketed intrastate. Prior to receiving approval for marketing a product, the manufacturer must prove its safety and efficacy through a rigorous testing program. Since the early 1950s, antibiotics have been used to stimulate growth and promote feed efficiency in ruminants, swine, and poultry. However, there is concern that continuous feeding of antibiotics to livestock may increase the resistance of certain strains of bacteria and that such resistance may be transferred to bacteria associated with human diseases. It is therefore expected that
the use of such antibiotics as penicillin and the tetracyclines in feeds may come under increased scrutiny. In 1954 the hormone diethylstilbesterol (DES) was approved for use in cattle finishing rations. Two years later, implants for steers were approved. In 1973 the use of DES was banned when it was proved to be carcinogenic and present in trace amounts in the meat of animals. However, in 1974 a Federal Court of Appeals overruled this ban. On November 1, 1979, the 1974 decision was reversed, and DES may no longer be used as a feed additive or implant for cattle and sheep. Another hormone, melengestrol acetate (MGA), is marketed for feeding to heifers being finished for slaughter. This compound, a synthetic progestin, affects metabolic rate and suppresses estrus, thereby increasing feed efficiency of feedlot heifers. Other hormones or hormonelike compounds which are marketed for implantation in beef cattle to improve growth and feed efficiency are Synovex and Ralgro. Synovex-S contains progesterone and estradiol benzoate, and is used for steers; Synovex-H contains testosterone and estradiol benzoate, and is used for heifers. Although both hormones are FDAapproved, they are being evaluated to determine the advisability of withdrawing approval of these products in consideration of human safety. Ralgro is the trade name for zeranol, a substance produced by the mold Giberella zeae. Although not a hormone itself, this compound is thought to influence the production and release of certain hormones in the body. It is approved for implantation in beef cattle. In fermenting feeds, rumen microorganisms produce carbon dioxide and methane in addition to fatty acids. Methane accounts for about 10% of the total energy of feedstuffs; hence its production reduces the energy efficiency of feeds. Therefore, compounds which may inhibit methane production in the rumen are being investigated. Dichlorovinyldimethyl phosphate, hemiacetal derivatives of chloral and starch, and halogenated analogs of methane have shown promise. A compound which has proved effective in altering rumen metabolism is monensin, a fermentation product of Streptomyces cinnamonensis. The proprietary product, Rumensin, has been shown to significantly improve feed efficiency in feedlot cattle. This compound acts by altering the proportion of fatty acids produced in the rumen; more propionic acid (and less methane) is produced than acetic and butyric acids. Because conversion of feed to propionic acid is energetically more efficient than its conversion to the other acids, overall production efficiency increases. See NUTRITION. International feed nomenclature. Because feed names have evolved through usage, the name of a particular feedstuff has not always been descriptive enough to fully identify and characterize it. For example, the term Linseed meal does not convey whether the product is obtained by mechanical or solvent extraction or whether it contains 33%, 35%, or 37% crude protein. To overcome this deficiency, a
691
692
Animal growth system was developed that makes it possible to identify the contents and other characteristics of a feed from its name. This system, known as the National Research Council (NRC) Feed Nomenclature System, or International Feed Nomenclature System, gives to each feed a specific name that may consist of up to nine terms: scientific name, origin (or parent material), species (variety or kind), part actually eaten, process(es) and treatment(s) to which the parent material or the part eaten was subject before being fed to the animal, stage of maturity (applicable only to forages), cutting or crop, grade (quality designations and guarantees), and classification (according to nutritional characteristics). In addition, each feed is grouped into one of eight feed classes, each of which is designated by a number in parenthesis: (1) dry forage or dry roughage; (2) pasture, range plant, and forages fed green; (3) silages; (4) energy feeds; (5) protein supplements; (6) minerals; (7) vitamins; (8) additives; then each feed is given an international reference number. The classification number, for example, (1), is the last term in the international feed name and the first digit of the feed reference number. An example of an international feed name for dehydrated alfalfa pellets is: alfalfa, aerial part, dehydrated ground pelleted, early bloom, cut 2, (1). The international reference number for this feed would be 1-07-733. This six-digit reference number is the numerical name of the particular feedstuff and may be used when formulating rations using linear programming and computers. Only the first digit is related to a specific feed characteristic. The remaining digits are used only to number the feed. This international system of naming feeds is being widely adopted for use in scientific writing and journal articles. International reference numbers particularly are being used to identify feeds in tables and computer formulations. Donald Burdick Pet Foods Pet foods are usually nutritionally complete products which adequately support growth and life, even if fed exclusively. Pet foods are made from proper mixtures of meat, fish, fowl, bone, cereal, milk, minerals, and vitamins. The minerals and vitamins are added in quantities necessary to augment the natural foodstuffs. They may be prepared in various finished forms, canned, frozen, or dried. The canned products are sterilized at high temperatures to permit storage without refrigeration. The frozen variety often is limited to raw meat or seafood. Dry items are stable at room temperature, though usually not as long as canned items. Those pet foods which include meat contain a fairly large proportion of highly nutritious organ meats as well as other meat products. The manufacture of the canned variety is fairly standardized in that the grinding, blending, filling, and sterilizing procedures are well known. The dry variety, particularly dry dog food, can be manufactured in a number of different ways, ranging from pelletizing of dry ingredients to hot-air and vacuum dry-
ing of moistened, expanded pastes and extrusions. Little or no preparation is needed for the frozen products since these are generally raw. A significant innovation is intermediate-moisture pet food that keeps without canning or refrigeration. Careful control of amount of “active water” in the product avoids spoilage. Aseptic canning of certain pet foods appears entirely possible. In this procedure the product is sterilized outside the can and then filled into a sterilized container, whereupon it requires no further cooking and is stable at room temperature for indefinite periods. Clarence K. Wiesman Bibliography. P. R. Cheeke, Applied Animal Nutrition: Feeds and Feeding, 2d ed., 1998; D. C. Church, Livestock Feeds and Feeding, 4th ed., 1997; A. Cullison and R. S. Lowrey, Feeds and Feeding, 4th ed., 1986; R. M. Ensminger, The Stockman’s Handbook, 7th ed., 1992; M. E. Ensminger, J. E. Oldfield, and W. W. Heinemann, Feeds and Nutrition, 2d ed., 1998; D. Natz (ed.), Feed Additive Compendium, 1977; W. G. Pond et al., Basic Animal Nutrition and Feeding, 1995.
Animal growth The increase in size or weight of an animal. A number of processes are involved in growth. These include hyperplasia, or increase in cell number due to cell proliferation (cell division or mitosis) and recruitment from stem cell populations; hypertrophy, or increase in the size of cells; and differentiation (precursor cells achieving mature functioning—often this process cannot be reversed). Schematic diagrams showing muscle and adipose development are included below. See ANIMAL MORPHOGENESIS; CELL DIFFERENTIATION; CELL DIVISION; CELL SENESCENCE AND DEATH; EMBRYONIC DIFFERENTIATION; MITOSIS. Measurement of growth. The indices of growth include height (as for humans and horses), body weight (agricultural animals), dry weight, length (rodents, lizards, fish), total body DNA or nitrogen retention (used in agricultural animals as an indicator of muscle growth), or rates of whole-body protein synthesis and/or accretion: Protein accretion (net protein synthesis) = Protein synthesis − Protein degradation Mathematical models for growth. Growth appears to follow an S or sigmoidal curve. There are multiple equations that mathematically describe growth. These include the logistic model and the Gompertz equation W = A · e−b · e−kt dW/dt = k · W · (A/W ) The Gompertz equation describes growth of an animal within a specific species very well from embryonic stage to mature adult size. Here W is the weight at any given age (time t following the beginning of
Animal growth development), A is the maximal weight (mature or adult weight), k is the growth rate or constant, and b is a second constant. At significant periods of growth, a linear function can easily be applied for growth. Telomeres. Telomeres are multiple TTAGGG tandem repeats at the end of chromosomes in at least higher vertebrates. The telomeres stablize the chromosomal ends and protect the genes in the subtelomeric regions. Telomeres shorten with progressive cell cycles (either in the animal or in tissue culture), leading to the view that shortened telomeres are related to senescence or aging. Addition of telomeric repeats can be accomplished by the enzyme telomerase, an RNA-dependent DNA polymerase. High levels of telomerase activity are observed in tissues with a high rate of cell proliferation, for example, malignant tissues and fetal tissues. Given the number of cell divisions, it is not surprising that cloned animals resulting from nuclear transfer from somatic cells have reduced telemetric length, with corresponding effects on growth rate. See CHROMOSOME. Embryonic and fetal development. Development begins following fertilization of the ovum (egg) by the spermatozoon. There is cell division to form first a two-cell stage, then four-cell, then eight-cell stage. Cell divisions continue but become desynchronized. Prior to implantation, in mammals, the part of the conceptus that will develop into the animal (as opposed to the membranes and placenta) is called the embryo. At the end of embryonic development, the major organs are normally formed and “true” bone formation begins (replacing cartilage). The fetus develops in the uterus until birth. Considerable growth and development occurs with organs achieving the capability of full functioning. In mammals, the placenta plays an important role influencing fetal growth. The availability of critical nutrients and oxygen to the growing fetus is related to placental size and functioning. In addition, the placenta can produce hormones that influence fetal growth. Other critically important hormones affecting fetal growth are insulin and insulinlike growth factor-II (IGF-II). In birds, the yolky egg provides the nutrients for the developing embryo. In amphibians, there is a free-living larval (for example, “tadpole”) stage that metamorphoses into the adult form as stimulated by the hormone thyroxine. See EGG (FOWL); HORMONE; PLACENTATION. Birth and postnatal growth and development. At birth, all animals have their organs formed. However, animals can be born (mammals) or hatched (birds) either with their systems well functioning (precocial development) or with the systems not completely functional (altricial development). For instance, horses can walk very soon after birth, where as cats and rats have offspring that are blind at birth and rats have hairless offspring. Following birth, the vast majority of growth (increase in weight) occurs. There is not complete synchrony of growth with different organ systems. First, the rate of skeletal growth is maximal, then comes maximal rate of muscular growth, followed
by maximal rate of adipose growth. Growth in most mammals ceases shortly after puberty (as discussed below). At puberty, there is an increase in sex steroid hormones in the circulation. The hormones induce a “growth spurt” in many species, including humans. Another result is the closure of the growth plate and the cessation of growth (discussed below). Muscle growth. Muscle development and growth involves the massive cell division (proliferation) of the myogenic precursor cells. These cells migrate to specific sites in the body, for instance, the developing limbs. They are attracted to the sites by gradients of growth factors and are undergoing some differentiation. The precursor cells then form myoblast cells, which are the last stage of cells that can undergo division. The myoblasts differentiate into myocytes. Multiple myocytes then aggregate to form multinucleated myotubes, which start to produce the contractile proteins and then fibers of muscle. Postnatally, muscle growth is accomplished by increasing the size of the muscle fibers. More nuclei are added by the recruitment of satellite cells. Muscle growth and development may be summarized schematically: Mesodermal somites ↓ Myogenic precursor cells (→ Satellite cells) ↓ Proliferating, differentiating, and migrating (embryonic and fetal life) ↓ Myoblast cells ↓ Proliferating and differentiating myocyte cells (capable of synthesizing many muscle proteins) ↓ Differentiating Myotubes (multinucleated cells that are assembling muscle contractile fibers) ↓ Myofibers (multinucleated) ↓ Muscles growing in size due to hypertrophy as well as recruitment of satellite cells See MUSCLE. Bone and cartilage growth. Longitudinal growth or lengthening of the long bones of the skeleton occurs at the epiphyseal or growth plate. Growth of this cartilage (including proliferation of cartilage cells or chondrocytes) occurs under the influence of the hormone insulinlike growth factor (IGF-I) (discussed below). At puberty or sexual maturation in most higher vertebrates, there is closing or ossification of epiphyseal plates, and longitudinal growth of bones ceases. This occurs under the influence of sex steroid hormones. See BONE; CARTILAGE. Adipose tissue growth. Adipose tissue provides energy storage as triglyceride. This is critically important to animals at a time of nutritional deprivation or during hibernation or when energy consumption is not meeting their needs (for example, during pregnancy and lactation). Unlike other organs, adipose
693
694
Animal growth tissue can continue to increase in size after sexual maturation. See ADIPOSE TISSUE; LIPID; TRIGLYCERIDE. Adipose tissue grows by proliferation and migration of adipose precursor cells. These differentiate into preadipocytes, which are the last cells in the adiocyte lineage capable of division. They may divide multiple times to form adipocytes. Adipocytes do not divide but produce the key enzymes necessary for triglyceride or fat formation and they fill with triglyceride. Adipose tissue growth and development may be summarized schematically: Adipose precursor cells ↓ Proliferating, differentiating, and migrating (embryonic and fetal life) ↓ Preadipocyte cells ↓ Proliferating and differentiating adipose cells (capable of synthesizing adipose proteins and triglyceride/fat) ↓ Mature fat-filled adipose cells (growing by increasing in size) Growth hormone. Multiple hormones influence growth. A prime example is growth hormone (GH), produced by specific cells, the somatotropes, found in the anterior pituitary gland. The development of the somatotropes is induced by pit-1. Release of GH is controlled by peptides released from specialized nerve terminals in the hypothalamus. These peptides include GH releasing hormone (GHRH) which stimulates GH release and somatostatin (SRIF) which inhibits GH release. Growth hormone acts by binding to GH receptors (GHR) on the cell membrane of cells in the liver and some other tissues. In turn, GH increases the release of IGF-I. The IGF-I binds reversibly to IGF-binding proteins (IGF-BPs) in the circulation. Unbound IGF-I stimulates skeletal growth: Hypothalamus (in brain) ↓ ↓ GHRH+ SRIF– ↓ ↓ Anterior pituitary gland (somatotropes) ↓ Growth hormone (GH) ↓ Liver and other tissues (GHR) ↓ ↓ IGF-I IGF-BPs ↓ ↓ Growing tissues, such as cartilage Lack of sufficient production of GH and/or IGFI is associated with dwarfism in, for example, humans, cattle, chickens, dogs, and mice. Excess GH production causes giantism and, in adult animals, acromegaly. Other hormones influencing growth. Normal levels of thyroid hormones are required for growth. Glucocorticoid hormones, produced by the adrenal cortex
[cortisol in humans and domestic mammals, and corticosterone in rodents and poultry and other birds], inhibit growth. Manipulation of growth rate. Some chemical agents (related to hormones) are used to stimulate growth in animals. These include the use of GH in children with abnormally low growth rates (and likely low adult heights) due to hypopituitary dwarfism or low height for age. Muscle growth is stimulated in many but not all higher vertebrates by androgens. In commercial cattle, growth can be enhanced by the administration of estrogens, alone or with androgens. This treatment is approved for commercial use in some countries. There are beta-adrenergic agonists (for instance, Ractopamine) that enhance growth, particularly muscle growth. This treatment is used commercially in pigs in some countries. Factors influencing growth rate. Growth rates within a species are influenced considerably by the environment and by the animal’s genetics. Growth may be reduced by lack of availability of nutrients, the presence of toxicants in the environment, environmental stressors (for example, high or low temperatures), and disease. Diets deficient in critical nutrients, for example, energy components, proteins (and even specific amino acids such as lysine, methionine, or tryptophan), vitamins, and minerals, will result in reduced growth rates. When full nutrition is restored, compensatory or catch-up growth occurs. Not only is the growth rate restored but it rebounds to a higher level, above that normally seen. One mechanism by which stressors reduce growth is via increased release of the glucocorticoid hormones, cortisol and corticosterone. Other mechanisms of growth reduction include suppression of IGF-I levels, for instance, by disease; and toxicants that influence the thyroid gland (goitrogens). See NUTRITION; PROTEIN. Genetics. Genetics profoundly influences growth rate. There is marked sexual dimorphism in growth rates, with males of most species growing faster and larger than females. Different breeds of animals have been produced with markedly different growth rates and mature body sizes. Perhaps the best example is seen in the widely different breeds of dogs and even the disparate-sized poodles such as miniature and toy poodles. Selective breeding can greatly influence growth rate. An example is seen in chickens (meat or broiler type), where since 1945 there has been increased growth rates by more than fivefold with consequent improvements in efficiency and reduced cost to the consumer. More recently, geneticists working for the livestock industry have employed quantitative trait loci (QTL) following genome scans or specific point mutations (for instance, single-nucleotide polymorphorisms, or SNPs) in candidate genes to identify superior parents or grandparents, and then the geneticists have used these in select breeding programs. See BREEDING (ANIMAL); GENETICS. Gene imprinting. Fetal growth rates are influenced by imprinting of genes. This is the result of methylation of DNA. The pattern of imprinting varies with the parent of origin of the genes. An example
Animal kingdom
PALEOZOIC
CHORDATA
CNIDARIA
CTENOPHORA
PLATYHELMINTHES NEMERTINES
MESOZOA GNATHOSTOMULIDA
NEMATODA
NEMATOMORPHA
KINORHYNCHA
LORICIFERA
GASTROPODA
GASTROTRICHA ROTIFERA
POLYPLACOPHORA
ONYCHOPHORA
TARDIGRADA
ACANTHOCEPHALA
CEPHALOPODA
MYRIAPODA CRUSTACEA
INSECTA
ANNELIDA
PHORONIDA PENTASTOMA
ECHIURA BRACHIOPODA ECTOPROCTA
PRIAPULIDA
VERTEBRATA
SIPUNCULA
TUNICATA CEPHALOCHORDATA
CHAETOGNATHA
POGONOPHORA ECHINODERMATA
MESOZOIC
CENOZOIC
One of five kingdoms of organisms: Animalia, Plantae, Fungi, Protoctista (or, in part, in some classifications Protozoa), and Procaryotae (bacteria). Animals are eukaryotic multicellular organisms that take food into their bodies and that develop from blastula embryos. More than 30 million species of organisms have been described in all five kingdoms; the animal kingdom includes over 1 million described species. Animal species are organized into phyla (see illus.) that are defined according to comparative patterns of development, body structures, behavior, biochem-
BIVALVIA
Animal kingdom
ENTOPROCTA
ical pathways, modes of nutrition, and ancestry. See ANIMAL SYSTEMATICS. Form and habitat. Animals take diverse forms. The largest animals are great blue whales, which grow up to 100 ft (30 m) in length; the smallest are marine loriciferans, measuring as little as 300 micrometers. Thus, the largest animals are about 100,000 times as large as the smallest. Some one-celled organisms are classified as members of the Protoctista kingdom. Most protoctists are microorganisms, but some are larger than animals. For example, rotifers, nematodes, and loriciferans are smaller than protoctists such as ciliates. Most animal phyla are aquatic; that is, they are found in shallow seas, ponds, puddles, and streams. Some members of two phyla are called terrestrial because they live on the land surface. These phyla are Chordata, including humans and other mammals, reptiles, amphibians, fish, and birds; and Arthropoda, including moths, beetles, and millipedes. Animals that fly are referred to as aerial; insects (phylum Arthropoda), birds, bats, and flying dinosaurs are included in this category. Some animals inhabit a moist, underground environment but are not strictly terrestrial; annelid worms, nematodes, and other soildwelling animals are adapted to life in the thin film of water surrounding soil particles. Diversity. There are approximately 1.2 million animal species, but new species and phyla are
ARACHNIDA
is the imprinting of the IGF-II gene in eutherian mammals. There are distinct differences in imprinting sites, depending on maternal or paternal origin of the gene. There is also imprinting of the IGF-II receptor gene (responsible for degradation of IGF-II); this is found in most mammals but not in the primate lineage. Colin G. Scanes Bibliography. S. Brody, Bioenergetics and Growth, Waverley Press, Baltimore, 1945; M. F. Rothschild, Porcine genomics delivers new tools and results: This little piggy did more than just go to market, Genet. Res., 83:1–6, 2004; C. G. Scanes, Biology of Growth of Domestic Animals, Iowa State Press/Blackwell, 2003.
ARTHROPODA PSEUDOCOELOMATES ACOELOMATES HEMI CHORDATA
MOLLUSCA
PRECAMBRIAN
PLACOZOA COELOMATES
PROTOCTIST ANCESTORS Animal phylogeny or family tree, diagramming the relationships among phyla of living animals. The ancestors of all present-day animals probably were soft-bodied marine forms, such as the placozoa. (After L. Margulis and K. V. Schwartz, “The Animal Kingdom” Teachers’ Guide, Wards Natural Science, Rochester, New York, 1987)
PORIFERA
695
696
Animal reproduction constantly being identified. Of the known animal species, some 75% are arthropods. One phylum, Vestimentifera, includes giant tube worms that measure 5 ft (1.5 m) long. These animals of the abyss live in hot, submarine gardens some 1.5 mi (2.5 km) below the ocean surface. Another phylum, Loricifera, includes spiny marine animals that live in offshore sediments. Much remains to be learned about loriciferan food and development because these animals stick tenaciously to shelly gravel. Unfortunately, animals along with their habitats are being obliterated thousands of times faster than any of the earlier mass extinctions recorded in the geological record. Biologists estimate that about 99.9% of all animal species that ever lived are now extinct. See EXTINCTION (BIOLOGY). Nutrition and development. Unlike plants, seaweeds, and some bacteria which photosynthesize (autotrophs), animals are fed by other sources (that is, they are heterotrophs). Animals obtain food by ingestive or absorptive nutrition. Most animals take food into their bodies (ingestive), but some aquatic animals absorb dissolved nutrients through their cell membranes either from symbiotic partners or from the surrounding water. The energy source is organic compounds. Animals may reproduce sexually or asexually or by both methods. Sexual reproduction involves two different reproductive cells, a female egg cell and a smaller male sperm cell. At fertilization these two cells join and produce a zygote that is destined to become a new organism. The animal zygote divides to form a blastula embryo that is a hollow, liquid-filled ball of cells. Blastulas are found only in the animal kingdom. Animal species also vary broadly in their mode of embryonic development. Embryos of some species (human, for example) develop directly into smaller versions of the adult, but in other species the embryos develop first into a larva, an immature stage preceding adulthood. For example, frog embryos develop into a tadpole, which is a larval stage, and then the tadpole dramatically reorganizes through the process of metamorphosis into the adult form. Adult animals produce sex cells that carry on the continuous life cycle of the species. Different patterns of development are one line of evidence utilized to deduce ancestry of a phylum. In contrast to sexual reproduction, asexual reproduction is by a single individual without fertilization of egg by sperm. New individuals are produced asexually by various processes, such as budding by sponges, detaching of body parts by annelids, parthenogenesis by rotifers, and pseudogamy by ticks. Asexual reproduction occurs in about half of the animal phyla. See BLASTULATION; EMBRYOGENESIS; REPRODUCTION (ANIMAL). Multicellularity. Animals are multicellular, but since all five kingdoms include certain multicellular species multicellularity alone does not define an animal. Animals, like plants, fungi, and protoctists, but unlike bacteria, are also eukaryotic. Eukaryotes have membrane-bound nuclei; more than one chromosome; and organelles (that is, “little organs”
within a cell, such as mitochondria and a nucleus). In contrast to plant cells, animal cells lack cell walls (which contain cellulose). Moreover, animal cells intercommunicate and support one another by means of elaborate junctions. In addition, the cells of animals form tissues which specialize in a wide array of functions; examples of such tissues are muscle, skin, and blood. The only exception is animals of phyla Porifera (sponges) and Placozoa, which have poorly defined tissues. See CELL (BIOLOGY); EUKARYOTAE; PLANT CELL; TISSUE. Invertebrates and vertebrates. Traditionally, animals have been grouped into invertebrates (without backbones) and vertebrates (with backbones). Vertebrates include mammals, amphibians, reptiles, birds, fish, all members of the phylum Chordata. However, a few chordates, such as tunicates, are invertebrates. Members of all other animal phyla, more than 98% of all animal species, are invertebrates. Although invertebrates lack backbones, they achieve physical support by structures ranging from delicate glass spicules, to tough rings and rods, to hydrostatic pressure. The number of animal species within a single invertebrate phylum varies considerably: Placozoa comprises a single species, but the great phylum Arthropoda comprises more than 800,000 species. If tropical species were better described, the arthropods might include as many as 10 million living insect species. See AMPHIBIA; AVES; CHONDRICHTHYES; CHORDATA; MAMMALIA; OSTEICHTHYES; REPTILIA. Animal ancestry. Microscopic evidence, fossils, geographic distribution, function (physiology), behavior, embryonic patterns, and form (morphology) are lines of evidence from which contemporary classifications are drawn, both within the animal kingdom and in linking it to other kingdoms. The common ancestors of living animals probably evolved over 700 million years ago. Most biologists agree that animals evolved from choanomastigotes, which are microscopic protozoa ancestral to the sponges. Examination of ribosomal ribonucleic acid sequences indicate that animals and fungi share an evolutionary history: their common ancestor may well be a single-celled organism similar to choanomastigotes. See ANIMAL EVOLUTION; PHYLOGENY; PLANT KINGDOM. Karlene V. Schwartz Bibliography. R. M. Alexander, The Chordates, 1975; M. F. Glaessner, The Dawn of Animal Life, 1984; E. D. Hanson, Origin and Early Evolution of Animals, 1977; R. Lewin, Thread of Life, 1991; L. Margulis and K. V. Schwartz, Five Kingdoms: An Illustrated Guide to the Phyla of Life on Earth, 3d ed., 1998; S. P. Parker, Synopsis and Classification of Living Organisms, 2 vols., 1982; E. E. Ruppert and R. D. Barnes, Invertebrate Zoology, 6th ed., 1994.
Animal reproduction The formation of new individuals, which may occur by asexual or sexual methods. In the asexual methods, which occur mainly among the lower animals, the offspring are derived from a single individual.
Animal reproduction Sexual methods are general throughout the animal kingdom, with offspring ordinarily derived from the paired union of special cells, the gametes, from two individuals. Basic to all processes of reproduction is the origin of the new individual from one or more living cells of the parent or parents. There has been no acceptable demonstration of the origin of a new living organism from other material than preexisting living cells. See PREBIOTIC ORGANIC SYNTHESIS. Asexual reproduction. Asexual processes of reproduction include binary fission, multiple fission, fragmentation, budding, and polyembryony. Among the protozoans and lower metazoans, these are common methods of reproduction. However, the last-mentioned process can occur in mammals, including humans. Binary fission involves an equal, or nearly equal, longitudinal or transverse splitting of the body of the parent into two parts, each of which grows to parental size and form. This method of reproduction occurs regularly among protozoans, in which it is essentially the process of cell division, with complete separation of the daughter cells. To a limited extent, binary fission may be observed among metazoans, such as sea anemones, as longitudinal fission and among planarians as transverse fission. Multiple fission, schizogony, or sporulation produces several new individuals from a single parent. It is common among the Sporozoa, such as the malarial parasite, which form cystlike structures containing many cells, each of which gives rise to a new individual. The cyst cells arise from a series of divisions of the nucleus, followed by cytoplasmic divisions of the original cell. See SPOROZOA. Fragmentation is a form of fission occurring in some metazoans, especially the Platyhelminthes, or flatworms; the Nemertinea, or ribbon worms; and the Annelida, or segmented worms. The parent worm breaks up into a number of parts, each of which regenerates missing structures to form a whole organism. It occurs also in certain starfish, as Linckia, in which single arms may pinch off and regenerate a complete animal. Budding is a form of asexual reproduction in which the new individual arises from a relatively small mass of cells that initially forms a growth or bud on the parental body. The bud may assume parental form either before separation from the body of the parent as in external budding, or afterward, as in internal budding. External budding is common among sponges, coelenterates (Fig. 1), bryozoans, flatworms, and tunicates. Among certain of the coelenterates, such as the colonial hydroid Obelia, buds give rise to medusae, or jellyfish, rather than to the parentaltype polyp. The medusae represent the sexual generation. They are free-swimming and of separate sexes, producing eggs and sperm, respectively. Upon fertilization of the eggs the asexual, polyp-type individual develops. Another example of asexual individuals budding sexual individuals is found in the cestodes, or tapeworms. Here, the head or scolex, by which the animal is attached to the host tissue, produces
testes bud ovary 25 mm
Fig. 1. Budding in Hydra.
a series of segments, termed proglottids, each of which is a sexual individual. This phenomenon is also known as strobilization. In some species of sponges, coelenterates, bryozoans, and tunicates, budding may occur without separation of the buds and thus lead to the formation of an organic colony, accompanied in some cases by specialization of parts for particular functions, such as in the Portuguese manof-war, Physalia. Internal budding occurs among fresh-water sponges and bryozoans. In the sponges the internal buds, termed gemmules, consist of groups of primitive cells surrounded by a dense capsule formed by the body wall. If the parent animal dies as a result of desiccation or low temperature, the cells of the gemmules can later be released and form new sponges. In the bryozoans the similarly functioning buds are known as statoblasts. Polyembryony is a form of asexual reproduction, occurring at an early developmental stage of a sexually produced embryo, in which two or more offspring are derived from a single egg. Examples are found scattered throughout the animal kingdom, including humans; in humans it is represented by identical twins, triplets, or quadruplets. In some flatworms, polyembryony is illustrated by the rediae, each of which in turn produces many young tadpolelike cercariae. A striking example of polyembryony is found among the insects in the hymenopteran Litomastix, which is a parasite on the egg of the moth Plusia. The embryo of this wasp subdivides so extensively that about 1500 individuals are formed. In mammals, polyembryony is a regular feature of the development of armadillos. The four offspring produced at a single birth in the nine-banded armadillo (Dasypus novemcinctus) are identical quadruplets. The quadrupling process occurs during the late blastocyst stage of development. In humans, twins occur in about 1.1% of births, triplets in 0.012%, and quadruplets in 0.00014%. About onethird of the twins are identical and have thus arisen by polyembryony, probably occurring in a late blastocyst stage at the onset of gastrulation. Sexual reproduction. Sexual reproduction in animals assumes various forms which may be classified under conjugation, autogamy, fertilization
697
698
Animal reproduction (syngamy), and parthenogenesis. Basically, the various processes all involve the occurrence of certain special nuclear changes, termed meiotic divisions, preliminary to the production of the new individual. See GAMETOGENESIS; MEIOSIS. Conjugation occurs principally among the ciliate protozoans, such as Paramecium, and involves a temporary union of two individuals during which each is “fertilized” by a micronucleus from the other. In this process the macronucleus of each conjugant breaks down, and the micronucleus undergoes two meiotic divisions to form four nuclei, of which three degenerate. The fourth nucleus divides again, and each conjugant transfers one of these micronuclei to its partner, where it fuses with the stationary micronucleus. The conjugants then separate and the fusion-micronucleus in each divides three times. Of the eight nuclei, four form macronuclei and three degenerate. The exconjugants then divide twice, along with mitotic division of the micronuclei, so that each of the four cells obtains one macro- and one micronucleus. The parental state is thus restored, and further reproduction occurs by simple fission. The ability to conjugate again is not attained until after a large number of divisions. Conjugation does not ordinarily occur within clones, which are organisms derived by mitotic division from a single individual. Strains that are capable of conjugating with one another are designated mating types, rather than males or females, because there are no correlated morphological distinctions between them, and moreover a particular species may have several mating types that are interfertile and thus represent more than two sexes. In autogamy the nuclear changes described for conjugation take place, but since there is no mating, there is no transfer of micronuclei. Instead, the prospective migratory micronucleus reunites with the stationary one. The process may be considered related to parthenogenesis. Fertilization, or syngamy, comprises a series of events in which two cells, the gametes, fuse and their nuclei, which had previously undergone meiotic divisions, fuse. In metazoans, the gametes are of two morphologically distinct types: spermatozoa, or microgametes, and eggs, also called ova or macrogametes. These types are produced by male and female animals, respectively, but in some cases both may be produced by a single, hermaphroditic individual. The nucleus of the spermatozoon has half the number of chromosomes characteristic of the ordinary (somatic) cells of the animal. The nucleus of the ripe egg in some animals, for instance, coelenterates and echinoderms, also has attained this haploid condition, but in most species of animals it is at an early stage of the meiotic divisions when ready for fertilization. In the latter situation, the meiotic divisions of the egg, characterized by formation of small, nonfunctional cells termed polar bodies, are completed after the sperm enters, whereupon the haploid egg nucleus fuses with the haploid sperm nucleus. Fertilization thus produces a zygote with the diploid chromosome number typical of the somatic cells of the species (23 pairs in humans) and this is main-
polar bodies (haploid) egg nucleus (diploid) (haploid)
(a)
(c)
sperm nucleus (haploid)
zygote nucleus (diploid)
(b)
(d)
cleavage nuclei (diploid)
Fig. 2. Stages (a–d ) of fertilization in animals. In different species the sperm enters at one or another of the stages between a and b, during which the two meiotic divisions of the egg occur.
tained during the ensuing cell divisions (Fig. 2). See FERTILIZATION. In many protozoans, such as the flagellate Polytoma, the fusing gametes are not visibly different from one another nor from the ordinary organism. They are termed isogametes. In others they may differ (anisogametes), usually in size as in Chlamydomonas braunii; in some cases (C. coccifera), one gamete is motile and the other not. Some metazoans are able to produce functional gametes while still in a larval condition. Reproduction of this type, termed neoteny, is exhibited in amphibians by the axolotl, and also in various insects. The Larvaceae among the tunicates are considered to represent persistent larvae reproducing neotenically. Hermaphroditism is rare among the vertebrates but is found among most groups of invertebrates. Despite the production of eggs and sperm by the same individual, most hermaphroditic animals reproduce by cross-fertilization. In some cases, this is due to differences in time of ripening of the eggs and sperm. In others, with internal fertilization, crossing is assured by differences in location of sperm and egg ducts. In ascidians there is a physiological block to the union of the spermatozoon with an egg from the same individual. The latter situation illustrates that the interaction of egg and sperm in fertilization is not only generally species-specific but may also be individual-specific. See HERMAPHRODITISM. Parthenogenesis is the development of the egg without fertilization by a spermatozoon. It is listed as a form of sexual reproduction because it involves development from a gamete. Rotifers, crustaceans,
Animal virus
Animal symmetry Animal symmetry relates the organization of parts in animal bodies to the geometrical design that each type suggests. The term asymmetrical applies to most sponges and some protozoa because the body lacks definite form or geometry, as it cannot be subdivided into like portions by one or more planes. Spherical symmetry is exhibited by some protozoa, such as the Heliozoia and Radiolaria. The body is spherical with its parts concentrically around, or radiating from, a central point. Radial symmetry is exemplified by the echinoderms and most coelenterates. The body is structurally a cylinder, tall or short, having a central axis named the longitudinal, anteroposterior, or oral-aboral axis (see illus.). Any plane through this axis divides the animal into like halves. Often several planes, from the axis outward, can divide the body into a number of like portions, or antimeres, five in most echinoderms. Ctenophores and many sea anemones and corals possess biradial symmetry, basically radial but with some parts arranged on one plane through the central axis. Most animals have bilateral, or two-sided, symmetry, in which a
RADIAL oral
radial plane median or saggital plane
aboral axis frontal plane
dorsal
posterior
mouth
anterior
and insects are the principal groups in which it occurs naturally. It has also been induced (artificial parthenogenesis) in species from all the major phyla by various kinds of chemical or physical treatment of the unfertilized egg. Even in mammals, several adult rabbits have reportedly been thus produced. See DEVELOPMENTAL BIOLOGY. The honeybee provides a classic example of natural parthenogenesis. The males, or drones, all develop from unfertilized eggs and are haploid. The females, or workers and queen, arise from fertilized eggs and are diploid. Since the queen is inseminated only once, during the nuptial flight, she stores the sperm during her egg-laying life, which is 5 years or more, and can evidently permit or prevent them from fertilizing the eggs that are laid. In certain animals periods of parthenogenesis may alternate with fertilization. Thus aphids produce parthenogenetic females during part of the year, but as winter approaches males appear, whereupon fertilization ensues. The fertilized eggs hatch out into females the following spring. Among certain of the Lepidoptera, like the gypsy moth, Lymantria, parthenogenesis occurs at times when there appears to be a scarcity of males, and the unfertilized eggs can form males or females. There is also sporadic occurrence of parthenogenesis in chickens and turkeys. Mostly the embryos are abnormal and die at an early stage, but in the turkeys many develop quite far, and a few have been reared to adults. Parthenogenetic development may also occur in eggs that are produced during a larval stage, as in the gallfly, Miastor. This condition is termed pedogenesis and is the parallel of neoteny. See ESTRUS; OOGENESIS; OVUM; SPERM CELL; SPERMATOGENESIS. Albert Tyler; Howard L. Hamilton
transverse plane or cross section
ventral
longitudinal axis
BILATERAL
Types of symmetry and the axes, planes, and regions in animal bodies. (After T. I. Storer et al., eds., General Zoology, 6th ed., McGraw-Hill, 1979)
median or sagittal plane divides the body into equivalent right and left halves, each a mirror image of the other. Tracy I. Storer
Animal virus A small infectious agent that is unable to replicate outside a living animal cell. Unlike other intracellular obligatory parasites (for example, chlamydiae and rickettsiae), they contain only one kind of nucleic acid, either deoxyribonucleic acid (DNA) or ribonucleic acid (RNA), but not both. They do not replicate by binary fission. Instead, they divert the host cell’s metabolism into synthesizing viral building blocks, which then self-assemble into new virus particles that are released into the environment. During the process of this synthesis, viruses utilize
200 nm
Fig. 1. Electron micrograph of negatively stained, purified bovine papilloma virions with icosahedral symmetry. (Courtesy of David Hill and John P. Sundberg)
699
700
Animal virus cellular metabolic energy, many cellular enzymes, and organelles which they themselves are unable to produce. For this reason they are incapable of sustaining an independent synthesis of their own components. Animal viruses are not susceptible to the action of antibiotics. The extracellular virus particle is called a virion, while the name virus is reserved for various phases of the intracellular development. See DEOXYRIBONUCLEIC ACID (DNA); RIBONUCLEIC ACID (RNA). Morphology. Virions are small, 20–300 nanometers in diameter, and pass through filters which retain most bacteria. Due to this property, viruses were called filterable agents of disease. However, large virions (for example, vaccinia, which is 300 nm in diameter) exceed in size some of the smaller bacteria. The major structural components of the virion are proteins and nucleic acid, but some virions also possess a lipid-containing membranous envelope. The protein molecules are arranged in a symmetrical shell, the capsid, around the DNA or RNA. The shell and the nucleic acid constitute the nucleocapsid. In electron micrographs of low resolution, virions appear to possess two basic shapes: spherical and cylindrical. High-resolution electron microscopy and x-ray diffraction studies of crystallized virions reveal that the “spherical” viruses are in fact polyhedral in their morphology, while the “cylindrical” virions display helical symmetry. The polyhedron most commonly encountered in virion structures is the icosahedron, in which the protein molecules are arranged on the surface of 20 equilateral triangles. Based on these morphological features, viruses are
400 nm
Fig. 2. Electron micrograph of a thin section of squirrel cells afflicted with squirrel fibroma. The brick-shaped pox virions contain a characteristic dumbbell-shaped nucleocapsid. (Courtesy of David Hill and John P. Sundberg)
classified as helical or icosahedral (Fig. 1). Certain groups of viruses do not exhibit any discernible features of symmetry and are classified as complex virions (Fig. 2). Further distinction is made betweenvirions containing RNA or DNA as their genomes and between those with naked or enveloped nucleocapsids. These assignments are shown in the table, in which the major groups of animal viruses are listed. Viral nucleic acid. The outer protein shell of the virion furnishes protection to the most important component, the viral genome, shielding it from destructive enzymes (ribonucleases or
Classification of animal viruses Family∗
Prototype
Nucleic acid
Structure
Poxviridae Parvoviridae Reoviridae
Vaccinia virus Adeno-associated Reovirus
Complex, enveloped Icosahedron, naked Icosahedron, naked
Rhabdoviridae
Vesicular stomatitis virus rabies
Double-stranded DNA Single-stranded DNA Double-stranded RNA, fragmented Single-stranded RNA negative strand
Herpesviridae
Herpes simplex
Double-stranded DNA
Adenoviridae Papovaviridae
Adenovirus Simian virus 40
Retroviridae Paramyxoviridae
Rous sarcoma Newcastle disease
Orthomyxoviridae
Influenza
Togaviridae
Alpha: Sindbis
Double-stranded DNA Double-stranded DNA, circular Single-stranded RNA Single-stranded RNA negative strand Single-stranded RNA negative strand, fragmented Single-stranded RNA
Coronaviridae Arenaviridae Picornaviridae Bunyaviridae
Flavi: Yellow fever Avian infectious bronchitis Lymphocytic choriomeningitis Polio Bunyamwera
Bullet-shaped, helical nucleocapsid, enveloped Icosahedron, enveloped Icosahedron, naked Icosahedron, naked Complex, enveloped Helical nucleocapsid, enveloped Helical nucleocapsid, enveloped Icosahedron, enveloped
Single-stranded RNA
Complex, enveloped
Single-stranded RNA negative strand Single-stranded RNA Single-stranded RNA negative strand, fragmented
Complex, enveloped Icosahedron, naked Complex, enveloped
∗After F. Fenner, The classification and nomenclature of viruses: Summary of results of meetings of the International Committee on Taxonomy & Viruses in Madrid, September 1975, Virology, 71:371, 1976.
Animal virus deoxyribonucleases). The viral genome carries information which specifies all viral structural and functional components required for the initiation and establishment of the infectious cycle and for the generation of new virions. This information is expressed in the alphabet of the genetic code (sequence of nucleotides) and may be contained in a double-stranded or single-stranded (parvoviruses) DNA, and doublestranded (reoviruses) or single-stranded RNA. The viral DNA may be linear or circular, and the viral RNA may be a single long chain or a number of shorter chains (fragmented genomes), each of which contains different genetic information. Furthermore, some RNA viruses have the genetic information expressed as a complementary nucleotide sequence. These are classified as negative-strand RNA viruses. Finally, the RNA tumor viruses have an intracellular DNA phase, during which the genetic information contained in the virion RNA is transcribed into a DNA and integrated into the host cell’s genome. The discovery of this process came as a surprise, since it was believed that the flow of genetic information was unidirectional from DNA to RNA to protein and could not take place in the opposite direction, from RNA to DNA. The transcription of RNA to DNA was termed reverse transcription, and the RNA tumor viruses are sometimes referred to as retroviruses. A classification based on the relationship between viral genomes and their messenger RNA (mRNA) is shown in Fig. 3. See GENETIC CODE; TUMOR VIRUSES. When introduced into a susceptible cell by either chemical or mechanical means, the naked viral nucleic acid is in most cases itself infectious. Two exceptions are the negative-strand RNA viruses and the RNA tumor viruses. In these cases the RNA has to be first transcribed and reverse-transcribed, respectively, into the proper form of genetic information before the infectious process can take place. This task is carried out by means of an enzyme which is contained in the protein shell of the virion nucleocapsid. − DNA Class VI Retro + RNA
+ RNA
− DNA
− RNA Class IV Picorna Toga Corona
Class II Parvo + DNA −
Class I Herpeto Papova Pox Adeno
+ mRNA
Class III + − RNA Reo (fragmented)
Class V − RNA Paramyxo Orthomyxo (fragmented) Rhabdo Arena Bunya (fragmented)
Fig. 3. Classification based on the relationship between the viral genome and its mRNA. (After D. Baltimore, Expression of animal virus genomes, Bact. Rev., 35:235–241, 1971)
The whole nucleocapsid is therefore required for infectivity. Assay. Advances in virus research are closely linked with the utilization of new laboratory techniques. The introduction of tissue cultures eliminated the necessity to use experimental animals and revolutionized animal virus research. The application of plaque assays, first developed in studies of bacterial viruses, was of similar importance. In this assay, a monolayer of cells growing on solid support (usually in a petri dish) is infected with a dilute solution of virions and overlaid with a soft nutrient agar layer. The high viscosity of the agar layer prevents newly released virions from diffusing farther than the cells adjacent to the ones originally infected. After a sufficiently large number of surrounding cells are killed by newly released virions (after 1 day to 3 weeks, depending on the virus), a clear plaque becomes visible. Since each plaque was initiated by a single virus particle, the method allows a quantitative determination of the number of viable virions and offers a means of virus purification (the virus population isolated from a single plaque is relatively homogeneous). Plaque assays are frequently used in diagnosis and in the clinical isolation of viruses. Other assays are based on the ability of some viruses to hemagglutinate (aggregate) red blood cells. The highest dilution of a virus-containing extract which is still capable of hemagglutinating (end-point dilution) is used as a measure of virion concentration. Titrations with specific antibodies are similarly applied in the characterization and quantitation of viruses. See TISSUE CULTURE. Infectious cycle. Viral infection is composed of several steps: adsorption, penetration, uncoating and eclipse, and maturation and release. Adsorption. Adsorption takes place on specific receptors in the membrane of an animal cell. The presence or absence of these receptors determines the tissue or species susceptibility to infection by a virus. Enveloped viruses exhibit surface spikes which are involved in adsorption; however, most animal viruses do not possess obvious attachment structures. Penetration. Penetration takes place through invagination and ingestion of the virion by the cell membrane (phagocytosis or viropexis). Enveloped viruses often also enter the cell by a process of fusion of the virion and cell membranes. Penetration is followed by uncoating of the nucleic acid, or in some cases by uncoating of the nucleocapsid. At this stage, the identity of the virion has disappeared, and viral infectivity cannot be recovered from disrupted cells (special methods are required to utilize the infectious nucleic acid). See PHAGOCYTOSIS. Eclipse. The absence of infectious particles in cell extracts is characteristic of the eclipse period. During the eclipse the biochemical processes of the cell are manipulated to synthesize viral proteins and nucleic acids. The survival of viruses depends on this subversive ability. In the process of evolution they have developed an extraordinary efficiency and a remarkable repertoire of strategies. The eclipse period in infections with DNA viruses starts with the
701
702
Animal virus transcription of the genetic information in the nucleus of the cell (except poxviruses), processing into mRNAs, and their translation into proteins (in the cytoplasm). This process is divided into early and late transcription. In the absence of newly synthesized viral components, the immediate early events must be entirely catalyzed by cellular enzymes. The early proteins, however, are virus-encoded functional proteins which will participate in the synthesis of viral DNA and of intermediate and late viral proteins, as well as in the shutoff of various cellular functions which might be detrimental to viral synthesis. The event that distinguishes early from late mRNA transcription and translation is the onset of viral DNA synthesis. The major late products are the structural proteins of the nucleocapsid. Almost as soon as these proteins are synthesized, they assemble with newly synthesized DNA molecules into virion nucleocapsids. The appearance of these particles signals the end of the eclipse period. The events of the eclipse period in infections with RNA viruses are similar, except that they take place in the cytoplasm (influenza virus excepted), and a division into early and late transcription cannot be made. In the case of positive-strand RNA viruses, the viral RNA is itself the mRNA. In infections with negativestrand RNA viruses, the virion RNA in the nucleocapsid is first transcribed into positive mRNAs. Intracellular nucleocapsids are present throughout the entire infectious cycle, and the eclipse period cannot be defined in the classical sense. RNA tumor viruses reverse-transcribe their RNA into DNA, which enters the cell nucleus and becomes integrated into the cellular DNA. All viral mRNAs and genomic RNAs are generated by transcription of the integrated DNA (Fig. 3). Maturation and release. The event characteristic of the maturation step is virion assembly and release. In many cases the protein shell is assembled first (procapsid) and the nucleic acid is inserted into it. During this insertion, processing of some shell proteins by cleavage takes place and is accompanied by a modification of the structure to accommodate the nucleic acid. Unenveloped viruses which mature in the cytoplasm (for example, poliovirus) often exit the cell rapidly by a reverse-phagocytosis process, even before the breakdown of the cell. In some cases, however, a large number of virus particles may accumulate inside the cell in crystalline arrays called inclusion bodies (Fig. 4). Viruses that mature in the nucleus are usually released slowly, and the damage to the cell is extensive. Enveloped viruses exit the cell by a process of budding. Viral envelope proteins (glycoproteins) become inserted at various sites into the cell membrane, where they also interact with matrix proteins and with nucleocapsids. The cellular membrane then curves around the complex and forms a bud which detaches from the rest of the cell. Figure 5 shows virus particles budding from the cell surface of a mouse cell infected with murine leukemia virus. Effect of viral infections. Two extreme types of effects are identified: lytic infections, which cause cell
400 nm
Fig. 4. Electron micrograph of thin section of chicken cells containing crystalline inclusion bodies formed by avian reovirus during infection. (Courtesy of David Hill and John P. Sundberg)
death by a variety of mechanisms with cell lysis as the most common outcome, and persistent infections, accompanied either by no apparent change in the host cell or by some interference with normal growth control, as in transformation of normal to cancer cells. The degenerative phenomena in tissue cultures during a lytic infection are called cytopathic or cytolytic effects. In animals, extensive destruction of tissue may accompany an infection by a lytic virus. See LYTIC INFECTION. As a defense to certain conditions of infection, animal cells generate a group of substances called interferons which, by a complex mechanism, inhibit replication of viruses. They are specific to the cell species from which they were derived but not to the virus which elicited their generation. (Mouse interferon will protect mouse but not human cells from any viral infection.) Genetics. Like other microorganisms, viruses mutate either spontaneously or following exposure to
1 µm
Fig. 5. Scanning electron micrograph of the surface of a mouse cell infected with murine leukemia virus. A large number of virus particles are shown in the process of budding. (Courtesy of R. MacLeod)
Animal virus mutagenic agents. With the introduction of cell tissue cultures and with the development of plaque purification procedures, isolation and characterization of well-defined mutants became possible. In the past the most useful mutants in genetic studies of animal viruses were temperature-sensitive (ts) and host-range (hr) mutants. The ts mutants grew efficiently at lower than normal but not at normal temperatures (permissive and nonpermissive temperatures, respectively). The mutation caused a change of an amino acid in a viral protein which resulted in a biologically inactive configuration at the higher temperatures. Studies of the impaired proteins and their respective genes helped to elucidate the molecular mechanisms of various events in the infectious cycle. The hr mutants, which grew normally in certain host cells but not in others, were useful in studies of virus–host interactions. This approach to viral genetics was dependent on the isolation of the proper mutants, which was often fortuitous. With the advent of modern DNA technology of restricting, cloning, and sequencing, it became possible to generate site-specific mutants. Rather than depending on an accidental change of a nucleotide by a mutagen, the investigator could chemically alter a nucleotide at any selected site. In this manner, mutants were generated to answer very specific questions about viral gene expression and virus–host interactions. Since reverse transcription can be carried out in the test tube, these techniques could, in principle, be applied also to RNA viruses as long as their RNAs are infectious. DNA cloning, which permits the preparation of relatively large quantities of viral genomes, is making it possible to study viruses which were inaccessible to investigation either because of low yields obtained in living organisms or because of the lack of suitable tissue cultures. See GENETIC ENGINEERING; MUTATION; RECOMBINATION (GENETICS); RESTRICTION ENZYME. This modern approach to animal virus genetics also rapidly advanced understanding of the complex molecular biology of the animal cell. Many structural and regulatory features of animal virus genes are similar to those of the animal cell, yet are easier to study because of the smaller size of viral genomes. See GENETICS; MOLECULAR BIOLOGY. Pathology. Virus infections spread in several ways: through aerosols and dust, by direct contact with carriers or their excretions, and by bites or stings of animal and insect vectors. At the point of entry, infected cells undergo viremia. From there, the virus becomes disseminated by secretions. It is carried through the lymphatic system and bloodstream to other target organs, where secondary viremias occur (except in localized infections like warts). In most cases viral infections are of short duration and great severity. However, persistent infections are not uncommon (herpes, adeno, various paramyxoviruses like measles). See VIRUS INFECTION, LATENT, PERSISTENT, SLOW. The afflicted organism mounts a variety of defenses, the most important of which is the immune response. Circulating antibodies against viral proteins are generated. Those interacting with virion
surface proteins neutralize the infectious potential of the virus. Although the antibodies are specific against the virus which has elicited them, they will cross-react with closely related virus strains. The specificity of neutralizing antibodies obtained from experimentally injected animals is utilized for diagnostic purposes or in quantitative assays. In addition to the circulating antibodies, cell-mediated immune responses also take place. The most important of these is the production of cytotoxic thymus-derived lymphocytes, found in the lymph nodes, spleen, and blood. They destroy all cells which harbor in their membrane viral glycoproteins, regardless of whether these cells are actively involved in virus synthesis or acquired the viral proteins by membrane fusion with inactive virions or cell debris. The cell-mediated immunity has been demonstrated to be more important to the process of recovery than circulating antibodies. In spite of their beneficial role, immune responses often seriously contribute to the pathology of the disease. Circulating antigen–antibody complexes can lodge in organs and cause inflammation; cell-mediated responses have been known to produce severe shock syndromes in patients with a history of previous exposures to the virus. See ANTIBODY; AUTOIMMUNITY; CELLULAR IMMUNOLOGY; IMMUNITY; VIRUS INTERFERENCE. Control. Viruses are resistant to the antibiotics commonly used against bacterial infections. The use of chemotherapeutic agents with antiviral activity is plagued by their toxicity to the animal host. However, the application of vaccines has been successful in the control of several viruses. The vaccines elicit immune responses and provide sometimes lifelong protection. Two types of vaccines have been applied: inactivated virus and live attentuated virus. Various inactivation procedures are available. For example, in preparations of polio vaccines (Salk), a combination of formalin and heat was utilized. Protection by attenuated live virus vaccines dates back as early as 1776, when Edward Jenner used infections with harmless cowpox to achieve protection from the devastating effects of smallpox. In modern times, an attentuated laboratory strain of smallpox has been applied so successfully that the disease is considered to be eradicated. A small probability of back mutations of the attenuated virus to a virulent strain makes applications of live vaccines somewhat riskier. On the other hand, protection is longer-lasting and, by virtue of spread to nonvaccinated individuals, more beneficial to the population group (herd effect). See CHEMOTHERAPY AND OTHER ANTINEOPLASTIC DRUGS; POLIOMYELITIS; VACCINATION; VIRUS CHEMOPROPHYLAXIS. In order to achieve full protection, it is important that the vaccine contain all the distinct antigenic types of the virus. Development of monoclonal antibodies led to a better characterization of these types in naturally occurring viruses. This information will undoubtedly lead to better vaccines. Moreover, monoclonal antibodies have aided investigations into the molecular structure of viral antigenic groups and brightened future prospects for synthetic vaccines.
703
704
Anise See MONOCLONAL ANTIBODIES; VIRUS; VIRUS CLASSIFICATION. M. E. Reichmann Bibliography. R. Dulbecco and H. S. Ginsberg, Virology, 1980; H. Fraenkel-Conrat and P. C. Kimball, Virology, 1982; H. Fraenkel-Conrat and R. R. Wagner (eds.), Comprehensive Virology, 1974–1983; R. A. Lerner, Synthetic vaccines, Sci. Amer., pp. 66–74, February 1983; S. E. Luria et al., General Virology, 3d ed., 1978; D. Nathans, Restriction endonucleases, Simian Virus 40 and the new genetics, Science, 206:903–909, 1979; K. M. Smith and M. A. Lauffer (eds.), Advances in Virus Research, 1958–1983.
Anise One of the earliest aromatics mentioned in literature. The plant, Pimpinella anisum (Apiaceae), is an annual herb about 2 ft (0.6 m) tall and a native of the Mediterranean region. It is cultivated extensively in Europe, Asia Minor, India, and parts of South America. The small fruits are used for flavoring cakes, curries, pastry, and candy. The distilled oil is used in medicine, soaps, perfumery, and cosmetics. See APIALES; SPICE AND FLAVORING. Perry D. Strausbaugh; Earl L. Core
Ankerite The carbonate mineral Ca(Fe,Mg)(CO3)2, also commonly containing some manganese. The mineral has hexagonal (rhombohedral) symmetry (see illus.) and has the cation-ordered structure of dolomite. The name is applied only to those species in which at least 20% of the magnesium positions are occupied by iron or manganese; species containing less iron are termed ferroan dolomites. The pure compound,
(a)
(b) Ankerite. (a) Specimen of the mineral (specimen from Department of Geology, Bryn Mawr College). (b) Crystal habits (after C. Klein, Manual of Mineralogy, 21st ed., John Wiley & Sons, Inc., 1993).
CaFe(CO3)2, has never been found in nature and has never been synthesized as an ordered compound. See DOLOMITE. Ankerite typically occurs in sedimentary rocks as a result of low-temperature metasomatism. It is commonly white to light brown, its specific gravity is about 3, and its hardness is about 4 on Mohs scale. See CARBONATE MINERALS; METASOMATISM. Alan M. Gaines Bibliography. L. L. Y. Chang, R. A. Howie and J. Zussman, Rock-Forming Minerals, vol. 58: Non-silicates: Sulphates, Carbonates, Phosphates, Halides, 1995; W. A. Deer, R. A. Howie, and J. Zussman, RockForming Minerals, vol. 1a: Orthosilicates 2d ed., 1997.
Annelida The phylum comprising the multisegmented, invertebrate wormlike animals, of which the most numerous are the marine bristle worms and the most familiar the terrestrial earthworms. The Annelida (meaning little annuli or rings) include the Polychaeta (meaning many setae); the earthworms and fresh-water worms, or Oligochaeta (meaning few setae); the marine and fresh-water Hirudinea or leeches; and two other marine classes having affinities with the Polychaeta: the Archiannelida (meaning primitive annelids) are small heteromorphic marine worms, and the Myzostomaria (meaning sucker mouths) are parasites of crinoid echinoderms. See HIRUDINEA; MYZOSTOMIDA; OLIGOCHAETA; POLYCHAETA. General Characteristics These five groups share few common characters and little resemblance except that most have a wormlike body. Typically they are bilaterally symmetrical, lack a skeleton, and have a short to long linear body divided into rings or segments, which are separated from one another by transverse walls or septa. The mouth is an anteroventral or anterior vent at the forward end of the alimentary tract, and the anus posterodorsal or posterior at the hind end of the gut. Metamerism. The linear series of segments, or metameres, from anterior to posterior ends constitute the annelid body. These segments may be similar throughout, resulting in an annulated cylinder, as in earthworms and Lumbrineris. More frequently the successive segments are dissimilar, resulting in regions modified for particular functions. Each segment may be simple or uniannular, corresponding to a metamere, or it may be divided or multiannulate. The total number of segments varies from five to several hundred, according to kind. Segments may have lateral fleshy outgrowths called parapodia (meaning side feet), armed with special secreted bristles or rods, called setae and acicula; they provide protection and aid in locomotion. Setae are lacking in Hirudinea and some polychaetes. The body is covered by a thin to thick epithelium which is never shed.
Annelida In some polychaetes they are highly developed, with lenses and retina (Fig. 3), and may be at the end of an erect stalk. Swarming stages of benthic annelids may acquire large modified eyes which are an attractant at sexual maturity. Some pelagic polychaetes have large eyes which make up most of the anterior end (Fig. 4). See EYE (INVERTEBRATE).
pharynx
cuticle prostomium eyes
epithelial cells lens rods
Fig. 1. Eteone (Phyllodocidae), everted pharynx, simple prostomial eyes, in dorsal view. (After O. Hartman, The littoral marine annelids of the Gulf of Mexico, Publ. Inst. Mar. Sci., 2:7–124, 1951)
Sense organs. Annelids have sense organs of many kinds. Most conspicuous are those on the anterior, or head end, and on parapodia. Eyes which function as photoreceptors may be variously developed, simple (Fig. 1) to complex, small to large, absent to numerous, and obscure to bizarre. Structurally they are simple light-absorbing pigment spots, usually paired on the head, but also on body segments or anal end. In some sabellids they are on the tentacular crown and may number in the hundreds, or they are on the pygidium (Fig. 2) and aid in backward movement.
peristomial eyes
pygidial eyes
Fig. 2. Fabricia (Sabellidae), peristomial and pygidial eyes, in ventral view. (After O. Hartman, Fabricinae featherduster polychaetous annelids in the Pacific, Pac. Sci., 5:379–391, 1951)
Fig. 3. Cross section of eye of Eunice (Eunicidae). (After O. Pflugfelder, Ueber den feineren Bau der Augen freilebender Polychaeten, Z. wiss. Zool., 142:540–586, 1932)
proboscis
Fig. 4. Anterior end of Torrea (Alciopidae), large compound eyes and everted proboscis, in dorsal view. (After E. Claparede, ` Les Annelides Chetopodes du Golfe de Naples, Mem. ´ Soc. Phys. Geneve, ` 20:1–225, 1870)
Feelers or tactoreceptors are frequently on the prostomium as antennae; on anterior segments as long filiform or thick fleshy structures called tentacles; and on parapodia as cirri, papillae, scales, or tactile hairs. Cilia in bands or clusters or in grooves occur in specific patterns. The most conspicuous receptors are the large caruncles (nuchal organs) of amphinomid polychaetes, as simple (Fig. 5) or complex (Fig. 6) fleshy, folded, paired organs surrounding the cephalic structures. The nuchal organs of some syllids and phyllodocids, as well as the ciliated ridges of errantiate and sedentary polychaetes, illustrate the diversity and importance of these structures. Proprioreceptors are sense organs which respond to stimuli originating within the animal, such
705
706
Annelida
prostomial antenna
nuchal organ
branchia Fig. 5. Simple nuchal organ or caruncle of Amphinome (Amphinomidae), in dorsal view. (After O. Hartman, The littoral marine annelids of the Gulf of Mexico, Publ. Inst. Mar. Sci., 2:7–124, 1951)
antenna eyes nuchal organ
branchia Fig. 6. Anterior end of Hermodice (Amphinomidae) with complex caruncle, in dorsal view. (After O. Hartman, The littoral marine annelids of the Gulf of Mexico, Publ. Inst. Mar. Sci., 2:7–124, 1951)
as hunger. Trophic reactions to sense food and repelling organs to deter enemies are well developed in many annelids. Nuchal organs are ciliated, eversible or stationary pouches or grooves located in the anterior or neck region. They relay messages of the immediate environment of the individual. When partly or wholly enclosed as pouches, they may contain small grains either of secreted or foreign origin to form statocysts; they function in righting movements of the body. Setae are part of the sensory system for they detect changes in the environment through their basal attachments. Shallow-water species often have short, strong, resistant setae, whereas abyssal species have long, slender, simple setae. Each organ is unique and well adapted to its role in the development, growth, protection, and reproduction of the species involved. Chromatophores are cell clusters which change their shape and size to conform to the shadows of the animal’s background, responding to changes in light intensities, and therefore are generally protective. They are well developed in translucent pelagic
larvae which exist at the surface of the sea; they screen damaging intensities of light from delicate tissues. Oligochaetes and hirudineans, which generally lack eyes or special light receptors, are sensitive to light changes through the surface epithelium. See CHROMATOPHORE. Luminescence. Light production is known in many annelids. Some oligochaetes shed luminous mucus through the mouth or anus. Many polychaetes have special photocytes which luminesce. Some Chaetopterus colonies shed bright mucus from the palpi and parapodia. Polycirrus among the terebellids is highly luminescent. Pelagic spawning stages of Odontosyllis (a syllid) and Tharyx (cirratulid) flash bright light when disturbed. Some scale worms (polynoids) have light-producing cells on the undersides of elytra. Flashes are rhythmic and may continue for a minute or more. Luminescence may lure other individuals of the species or aid escape. The posterior, broken-off end continues to light up long after the anterior end has escaped. See BIOLUMINESCENCE. Digestive system. The alimentary tract is a straight or sinuous tube, consisting of mouth, pharynx, esophagus, stomach, gut, and pygidium or anus. Its several parts may be further differentiated and named for structure and function. The mouth may be a simple anterior (oligochaetes) or anteroventral (many polychaetes) pore provided with highly complex organs or accessory parts. The upper lip is formed by the prostomium, and the lower lip by the ventrum of the first segment. Accessory organs include the short to long grooved palpi (Fig. 7) of many polychaetes, which direct food to the mouth or also select and propel nutrients along ciliated tracts. The building organs and cementing glands of some tubicolous annelids are associated with the mouth; they select inert particles for shape and size and attach them in specific patterns to a basic secreted mucoid membrane, resulting in tubes which are highly characteristic (Fig. 8). In sand-cementing or reef-building colonial species, the beach front may be covered with cemented bands of tubes highly
pharynx palp
Fig. 7. Laetmonice (Aphroditidae), everted pharynx, long palpi, in dorsal view. (After W. McIntosh, A monograph of the British annelids, Roy. Soc. London, 1(pt. 2):215–442, 1900)
Annelida
(c)
(a) (b)
Fig. 8. Tubes of sedentary polychaetes. (a) Distal end of tube of Lanice (Terebellidae) (after O. Hartman, Polychaetous annelids collected by the USNS Eltanin and Staten Island cruises, chiefly from Antarctic seas, Allan Hancock Monogr. Mar. Biol., 2:1–387, 1967). (b) Tube of Vermiliopsis (Serpulidae), entire and in frontal view. (c) Several tubes of Protula (Serpulidae) on deep-water pelecypod (after O. Hartman, Endemism in the North Pacific Ocean, with emphasis on the distribution of marine annelids, and descriptions of new or little known species, Essays in the Natural Sciences in Honor of Captain Allan Hancock, University of Southern California, Los Angeles, 1955).
resistant to wave action. Colonial serpulid worms secrete calcium carbonate from special collar cells, formed into highly organized tubes with the aid of oral organs. The mouth is followed by the buccal cavity which may be modified as a proboscis or saclike eversible pouch. In nonselective, mud-eating annelids the pouch may be voluminous when everted, its surface papillated or branched so as to increase its sweeping surface as it brushes over the surface muds; the adherent materials are swept into the oral cavity and swallowed. The buccal cavity may be followed by a short to long muscularized eversible pharynx (Fig. 9) which captures and breaks up or compresses food particles. Its inner walls may be fortified by papillae (Fig. 10) or hard gnaths or jaws (Fig. 11) called para-, micro-, or macrognaths, maxillae, mandibles, chevrons, trepan, or other descriptive names. A short esophagus leads to the muscular stomach or digestive region. Lateral ceca or pouches may be present, along esophageal and stomach portions, to increase the amount of surface for secretion and digestion, especially in short-bodied worms. Peristaltic (clasping and compressing) movements are rhythmic and result in the food bolus being digested, and wastes separated and pushed into the gut. In nonselectivefeeding annelids the gut may be distended with great amounts of inert materials; in selective feeders there are few remains but those of living animals. The proctodaeum, or region preceding the anus, expels the wastes as fecal pellets of characteristic form. The proctodaeum with pygidium may be a simple, tapering end with posterior pore or may be highly modified as a plaque. In some pelagic swarming nereids the pygidium of the male acquires a circlet of papillae through which sperm cells are shed into the sea. Defecation, or disposal of wastes from
the gut, is usually posteriorly. Tubicolous annelids which lie head down may reverse their position for this function; others unable to turn may have an anal ciliated groove directing fecal materials forward and dorsal, away from the mouth. Nervous system. The nervous system consists of a dorsal, bilaterally symmetrical, ganglionic mass or brain within or behind the preoral region. The brain is connected to the ventral cord through the circumesophageal connectives which extend about the
pharynx
Fig. 9. Anterior end of Glycinde (Goniadidae), pharynx partly everted, in right lateral view. (After O. Hartman, Goniadidae, Glyceridae, Nephtyidae, Allan Hancock Pacific Expedition, 15(1):1–181, 1950)
pharynx Fig. 10. Anterior end of Anaitides (Phyllodocidae) with pharynx fully everted, in dorsal view. (After E. Rioja, Datos para el conocimiento de la fauna de Anelidos ´ poliquetos del Cantabrico, ´ Trab. Mus. Nac. Madrid, Zool. 37, pp. 1–99, 1918)
707
708
Annelida paragnaths
(a)
(b)
Fig. 11. Anterior end of Nereis (Nereidae) with everted pharynx showing distal jaws and paragnaths, (a) dorsal and (b) ventral views. (After E. Rioja, Datos para el conocimiento de la fauna de Anelidos poliquetos del Cantabrico, ´ Trab. Mus. Nac. Madrid, Zool. 37, pp. 1–99, 1918)
oral cavity (Fig. 12). The ventral cord may be single or paired, nearly smooth or nodular, or ganglionated according to the segmental pattern of the body. The brain sends out lateral branches to the eyes, palpi, antennae, or other structures; the ventral cord has lateral branches to all fleshy parts which receive stimuli. A giant axon is an enlarged part of the ventral cord, present in many long, muscular, or actively moving oligochaetes and polychaetes. It permits rapid transfer of stimuli and muscular response, resulting in abrupt response. Messages are relayed directly from the source of stimulation to the muscles, without reaching the brain. These giant fibers were first discovered in oligochaetes in 1861, and have since been found in many other annelids. They vary in size and complexity, but are always larger than normal neurons to which they are connected. Their cells arise in the circumesophageal ganglion and extend throughout the length of the worm, or they arise much farther back and may extend longitudinally or crosswise. Each giant fiber may represent a single nerve cell or fusion of several nerve cells. Their diameters
antenna brain
proboscis neuron Fig. 12. Diagram of nervous system of Nereis showing some of the nerve branches. (After S. Maxwell, Beitra¨ zur Gehirnphysiologie der Anneliden, Pfluegers Arch., Berlin, 67:263–297, 1897)
are largest in Myxicola (Sabellidae) and much smaller in most other annelids. Their occurrence has been reported from widely related families. Neurosecretion is a process of special cells located in the brain. Typically, nerve cells are polarized, consisting of an axon or head and a slender fiber. Secretory cells lack a fiber and secrete hormones. In some annelids they affect the time and rate of metamorphosis at maturity. In nereids the hormones inhibit the processes of metamorphosis, until the sexual cells have arrived at their full stage of development. The two processes, sexual development and transformation to an epitoke, are separate and each under its own control. See NEUROSECRETION. Circulatory system. The circulatory system consists of dorsal and ventral longitudinal, median vessels located above and below the alimentary tract. Lateral branches extend to all parts of the body. Pulsating or propelling contractile portions, sometimes called hearts, are in the anterior dorsal vessel or also at intervals in the ventral vessels. In oligochaetes these are segmental vessels of varying number connecting the dorsal and ventral vessels and surrounding some portion of the alimentary tract. The contained blood is red through the presence of a hemoglobinlike substance, which may be in corpuscles or granular in the fluid, or it is green, through the presence of chlorocruorin. These colors when diluted are yellow or colorless, as in many small annelids. See RESPIRATORY PIGMENTS (INVERTEBRATE). The purpose of the blood is to transfer nutrients to all parts of the body and to exchange carbon dioxide for oxygen, either through the skin or through special organs called branchiae (Fig. 5). They are fingerlike fleshy prolongations of the epithelium, either simple, branched, or spiraled structures, at anterior or also posterior ends of the body. Each branchia contains a vascular loop through which the blood flows in one direction. The oxygen requirements of annelids may be very low, illustrated in worms existing in stagnant pools and in deep oceanic basins. Irrigation in the tube, burrow, or environment may be achieved through U-shaped structures which open to the surface or through undulations of the body or parapodia, to ensure a flow of water and nutrients. Some annelids lack a closed circulatory system so that blood and coelomic fluids mix freely, resulting in a hemocoel; it may be partial or complete. Many annelids have a special organ called a cardiac or heart body surrounding the pulsating vessel and sometimes visible as a thick brown or red body of spongy tissue; its function is to dispose of circulatory wastes. The circulatory and coelomic fluids aid in maintaining turgidity of the annelid body, and with musculature control they act as a kind of skeleton. Excretory system. The excretory system functions to dispose of wastes in the coelom, resulting from metabolism. The special organs for the function are called nephridia. In their simplest form they are protonephridia, consisting of a strand of cells connecting the coelom to the body wall, usually a pair to a segment. The inner end terminates in a flagellum which propels wastes outward; the distal end is a
Annelida pore at the surface of the body called a nephridiopore. Primitive nephridia occur in larval, small, or primitive annelids. Clusters of simple nephridia, solenocytes, basally joined to a tube, are present in many errantiate polychaetes. More complex organs, or metanephridia, have a ciliated nephrostome or funnel opening into the coelom and continued to the surface as a complex organ. They function to transport wastes and at sexual maturity may serve as gonoducts to release gonadial products. Complex nephridia are present in all oligochaetes and many polychaetes. Nephridia may be few in number and conspicuous, or limited to anterior segments as in many sedentary families, or may be segmental, a pair to a segment, or dispersed and irregular. Some coelomic wastes, which may result through phagocytosis, or cell breakdown, cannot be eliminated through nephridia and may accumulate to form conspicuous dark blobs in posterior segments. Muscular system. The muscular system consists of an outer circular and an inner longitudinal system of muscles, each varying in extent and density according to species. In addition, an oblique series, between outer and inner layers, is well developed in annelids performing complex lateral movements. Long, very active burrowers or crawlers have an extensive musculature, whereas short, sluggish forms may have diminished musculature development. The circular muscles form a sheath beneath the epithelium; they effect an increase or decrease in body diameter, and are best developed in tubicolous and burrowing forms. Longitudinal muscles are often in paired thin to thick bands in dorsal and ventral positions; they effect a shortening or lengthening of the body, and are most conspicuous in errantiate species. The oblique muscles occur in incomplete bands, between circular and oblique layers, connecting parts of the parapodia to the body wall and to other muscular layers. Movements are achieved mainly by coordinated muscular contractions and expansions of the laterally projecting parapodia or setae, resulting in an undulating or meandering movement. Swimming species may move from side to side or by successive forward darts and stops. Some annelids with reduced parapodia and a long proboscis use the latter in progression by extension and withdrawal of the eversible part of the alimentary tract. Some annelids secrete mucus to move over irregular surfaces. Regeneration. The ability to replace lost parts is highly developed in annelids. Most frequent is the replacement of tail, parapodia, and setae. The anterior end may be replaced provided the break is postpharyngeal. The torn end is first covered over with scar tissue, then differentiated into epithelial cells and all other tissues characteristic of the whole animal. The tendency is to restore the number of segments peculiar to the species. Least differentiated species regenerate more completely than highly modified ones. Younger individuals repair more easily or quickly than older ones. Fragmentation may result in as many new individuals as there are pieces.
apical tuft eye prototroch
mouth mesodermal band
alimentary tract
anal tuft Fig. 13. Lateral view of larval trochophore with parts indicated and much enlarged.
The presence of a part of the ventral nerve cord is necessary to develop a new head. See REGENERATIVE BIOLOGY. Reproduction. Depending on the species, reproduction may be sexual, asexual, or both. Sexual reproduction may be dioecious, in which male and female are similar, rarely dissimilar. Individuals may be hermaphroditic, both male and female, but with cross fertilization. Some annelids are protandric hermaphrodites, in which the sexual stages alternate. Epigamy refers to the transformation of mature individuals from a benthic to a planktonic form; changes occur in parapodial lobes and the alimentary tract may be absorbed. Fully metamorphosed individuals leave their tubes, swarm to the surface of the sea, sometimes in rhythmic or lunar patterns, deposit sperm and ova, and then disintegrate. In the palolo (Eunicidae) only posterior ends are released, and head ends retreat to regenerate new tails. The ova which are deposited at the surface of the sea develop into ciliated spheres called trochophores (Fig. 13), which develop into segmented metatrochs (Fig. 14) and then into larvae that begin to resemble adults (Fig. 15) into which they will metamorphose.
Fig. 14. Metatroch of planktonic sabellarian larva, showing early segmentation and long swimming setae. (After O. Hartman, Polychaetous annelids, Paraonidae, Magelonidae, Longosomidae, Ctenodrilidae and Sabellariidae, Allan Hancock Pacific Expedition, 10(3):311–389, 1944)
709
710
Annelida
Fig. 15. Transitional stage of metatroch larva of sabellarian larva, with larval setae nearly replaced. (After O. Hartman, Polychaetous annelids, Paraonidae, Magelonidae, Longosomidae, Ctenodrilidae and Sabellariidae, Allan Hancock Pacific Expedition, 10(3):311–389, 1944)
At the appropriate time they sink to the bottom and assume a benthic development. See PROTANDRY. Olga Hartman Fossils Fossil annelids, or segmented worms, are mostly softbodied, yet sufficiently common to indicate that this large and varied group of invertebrates has been abundant for more than 500 million years. Origins. The bulk of the annelidan fossil record is represented by the polychaetes. The earliest definite occurrences are from the Lower Cambrian of Greenland (Sirius Passet fauna). Possible precursors to the polychaetes exist in the form of wiwaxiids and the more primitive halkieriids. Both are sclerite-bearing, and in the wiwaxiids the microstructure is identical to annelidan chaetae; in these animals the sclerites form a protective coat, but it is plausible to derive the notochaetae and neurochaetae from, respectively, an imbricated dorsal array and fanlike extensions of leglike parapodia. The coexistence of polychaetes and halkieriids in the Sirius Passet fauna indicates an earlier divergence, but a number of Ediacaran fossils previously assigned to the annelids, such as Dickinsonia, probably are not closely related. Soft-bodied polychaetes. The soft-bodied record from the Cambrian onward is sparse, but polychaetes are known from a number of exceptionally preserved fossil biotas. A diverse assemblage occurs in the Middle Cambrian Burgess Shale of British Columbia. Not only do these polychaetes (Fig. 16) show a very wide morphological range, but with one possible exception it has not proved feasible to accommodate any of the six species in a known fossil or living family. See BURGESS SHALE. Other landmarks in the polychaete fossil record are a series of deposits from the Carboniferous that include the Francis Creek Shale (Mazon Creek ironstone nodules; Fig. 17) of Illinois, the Bear Gulch Limestone of Montana, and the Granton shrimp bed of Edinburgh, Scotland. The first two units have yielded a wide variety of polychaetes, many with jaws. Although only one polychaete occurs at
Granton, it is significant in apparently representing the earliest tomopterid. In younger deposits, softbodied polychaetes are known from horizons in the Triassic (Voltzia Sandstone of eastern France, also deposits in Italy and Madagascar), Jurassic (Osteno beds of Italy, Solnhofen Limestone of Germany), Cretaceous (Sahel Alma fish beds of Lebanon), and Eocene (Monte Bolca of Italy). Scolecodonts. Polychaete jaws, or scolecodonts (Fig. 18), first appear in the lower Ordovician but seem to be conspicuously more abundant in the Paleozoic than in either the Mesozoic or Cenozoic. Many are attributed to the eunicids, the jaws of which seem to have a high preservation potential in comparison to those of many other families. Such variation is due to differences in original scolecodont mineralogy, which are still imperfectly understood. However, it seems clear that composition varies taxonomically, with different families having scolecodonts either composed of organic material (scleroproteins) or having this organic matrix mineralized to varying degrees by either aragonite or calcite. In Recent polychaetes, for example, eunicids and onuphids have calcite mineralization and lumbrinerids have aragonite, while arabellids and glycerids are
5 mm Fig. 16. Canadia spinosa from the Middle Cambrian Burgess Shale of British Columbia. Note the pair of anterior tentacles, broad dorsal chaetae extending across much of the trunk, and ventral chaetae extending from either margin.
Annelida unmineralized. It also appears that the composition of scolecodonts has changed during time. Paleozoic scolecodonts were predominantly organic and strongly sclerotized, whereas calcareous mineralization became widespread in the Mesozoic. This mineralogical shift may explain the scarcity of postPaleozoic scolecodonts, given that aragonite is relatively unstable, and most scolecodonts are extracted in the laboratory by acid digestion. Work on preservation potential, including survival during digestion by predators and resistance to decay in the sediment, is improving understanding, but the correlation between expected survival according to composition and field observation is by no means exact. Other fossil records. Apart from scolecodonts, the hard-part record of polychaetes consists of tubes, especially the calcareous serpulids and related spirorbids. Spirorbids first appear in the Ordovician, and both groups are a frequent component of many Mesozoic and Cenozoic assemblages. They may occur as encrusters on reefs, hard grounds, and shell debris, or as free-living on sandy substrates. Numerous groups of polychaetes inhabit tubes that are variously reinforced so as to give a wide spectrum of fossilization potential. Most important in the fossil record are the terebellids, sabellariids, and pectinariids, all of which construct an agglutinated tube bound together by mucous secretions. This may be of mud, sand grains, or other selected particles such as foraminiferal tests, bivalve shells, echinoderm ossicles, and fish scales. Assignment to particular groups largely depends on comparisons with living relatives. However, a remarkable specimen from the Miocene of New Zealand shows a soft-bodied sabellid within its tube. In the Paleozoic a variety of tubicolous fossils have been attributed to polychaetes, but many of these are of questionable status. It has also been customary to attribute many burrows and other trace fossils to polychaete activities, but in few cases are such assignments secure because a variety of vermiform phyla are capable of making effectively indistinguishable traces, but trace fossils (Walcottia) from the Upper Ordovician of Cincinnati are so similar to polychaetes that an assignment to this group seems secure. Characteristic borings into hard substrates, such as by spionid polychaetes, are also diagnostic. A unique specimen from the Middle Devonian of New York has a spionidlike polychaete still preserved in its boring. See TRACE FOSSILS. Oligochaetes and leeches. The fossil record of oligochaetes is lamentable, in part because their predominantly terrestrial and fresh-water habitats have a relatively poor rock record. Nevertheless, the reproductive cocoons characteristic of the clitellates (oligochaetes and leeches) have been recognized as far back as the Triassic and may be more common as fossils than is generally realized. In addition, there is a record of oligochaete trace fossils in the form of burrows and associated fecal pellets from ancient soils (paleosols). The most convincing of these come from Cenozoic paleosols, especially in the western interior of the United States. However, the oldest records appear to date back to the Triassic.
(a)
10 mm
(b)
711
10 mm
Fig. 17. Polychaetes from the Francis Creek Shale (Mazon Creek ironstone nodules). (a) Astreptoscolex anasillosus, a nephytid; note the marginal chaetae and cuticular creases. (b) Esconites zelus, a glycerid; the body is comparatively poorly preserved, but the anterior scolecodont apparatus (arrow) is in place.
Fig. 18. Scolecodont element (Kettnerites martinssoni) from the Silurian of Gotland, Sweden. (Courtesy of Claes Bergman, Kristianstad, Sweden)
The leeches are best known for their bloodsucking abilities, but they also include a number of predators. Apart from the cocoons, they are poorly represented in the fossil record. A tentative record from the Lower Silurian of Wisconsin is based on the putative terminal sucker of an annulated soft-bodied worm. Much younger worms from the Upper Jurassic Solnhofen Limestone are also claimed, on slender evidence, to represent leeches. Echiurans. Molecular data indicate that the echiurans, long thought to be related to the Annelida, may nest within the Polychaeta. The fossil record, however, is restricted to examples from the upper Carboniferous (Mazon Creek) and, more speculatively, various traces attributed to the activity of the characteristic proboscis.
712
Anomopoda Pogonophorans. Pogonophorans were once regarded as a separate phylum, but evidence (including molecular) now indicates that they nest within the annelids. Soft-bodied remains are not known, but certain tubicolous fossils have been assigned with varying degrees of confidence to the pogonophorans. Some are associated with hydrothermal deposits, ancient analogues of the midocean ridge “black smokers.” See FOSSIL. S. Conway Morris Bibliography. D. E. G. Briggs and E. N. K. Clarkson, The first tomopterid, a polychaete from the Carboniferous of Scotland, Lethaia, 20:257–262, 1987; G. K. Colbath, Jaw mineralogy in eucinean polychaetes (Annelida), Micropaleontology, 32:186–189, 1986; S. Conway Morris, Middle Cambrian polychaetes from the Burgess Shale of British Columbia, Phil. Trans. Roy. Soc. London, B285:227–274, 1979; S. Conway Morris and J. S. Peel, Articulated halkieriids from the Lower Cambrian of North Greenland and their role in early protostome evolution, Phil. Trans. Roy. Soc. London, B347:305–358, 1995; D. Jones and I. Thomson, Echiura from the Pennsylvanian Essex fauna of northern Illinois, Lethaia, 10:317– 325, 1977; S. B. Manum, M. N. Bose, and R. T. Sawyer, Clitellate cocoons in freshwater deposits since the Triassic, Zool. Scripta, 20:347–366, 1991; N. Kotake, Deep-sea echiurans: Possible producers of Zoophycos, Lethaia, 25:311–316, 1992; I. Thompson, Errant polychaetes (Annelida) from the Pennsylvanian Essex fauna of northern Illinois, Palaeontographica, Abteilung A, 163:159–199, 1979.
Anomopoda An order of fresh-water branchiopod crustaceans, formerly included in the Cladocera. Exceptionally these organisms may be as much as 6 mm (0.24 in.) in length, but often are less than 1 mm (0.04 in.). So-called water fleas of the genus Daphnia are the most familiar anomopods. Daphnia and its close relatives swim freely in open water; however, many anomopods are benthic, predominantly crawling species. The short trunk is protected by a functionally bivalved carapace that covers the five or six pairs of trunk limbs and the terminal postabdomen that bears two claws. The trunk limbs are extremely diverse in structure, both among themselves and between species. They collect food, usually particulate, by scraping or filtration or both, and are used by many benthic species for crawling. The postabdomen is often used for pushing. Swimming is by means of antennae. The short antennules (first antennae) are sensory. There is a single median sessile eye, derived by fusion of two eyes. Most species feed on minute particles, but Pseudochydorus is a scavenger, and Anchistropus has the remarkable habit of living parasitically on Hydra, against whose stinging cells it is protected. The mandibles are of the rolling, crushing type. Reproduction is mainly by parthenogenesis, eggs and young being carried dorsally in a brood pouch
beneath the carapace. Some races of some species of Daphnia are obligate parthenogens. Cyclical parthenogenesis is common in temperate zone species; summer parthenogenesis alternates with sexual reproduction in autumn. Fertilized eggs are cast off, protected by a modified portion of the carapace, which is often called an ephippium, and the eggs can overwinter or resist drought. Worldwide in distribution, anomopods are among the most abundant and successful of all fresh-water invertebrates. See BRANCHIOPODA. Geoffrey Fryer Bibliography. S. P. Parker (ed.), Synopsis and Classification of Living Organisms, 2 vols., 1982.
Anopla A class of the phylum Rhynchocoela which is divided into the orders Palaeonemertini and Heteronemertini. The simple tubular proboscis lacks stylets and resembles the body wall in structure. The mouth is posterior to the brain. The nervous system lies either immediately below the epidermis or among the musculature of the body wall. The vascular system is well developed. Anoplan rhynchocoelans are common on rocky shores, living beneath stones or interstitially in shell debris or coarser sands. They are carnivorous, preying almost entirely on annelids, but will also take carrion if this is not too decomposed. Acid secretions kill or immobilize prey in the foregut, and digestion occurs rapidly in the intestine by extra- and intracellular processes. See ENOPLA; HETERONEMERTINI; PALAEONEMERTINI; NEMERTEA. J. B. Jennings
Anoplura A small group of insects usually considered to constitute an order. They are commonly known as the sucking lice. All are parasites living in the covering hair of mammals. About 250 species are now known, and these comprise probably about half of the species in the world, but the group is known in detail to few persons. Morphology. These lice are distinguishable from the Mallophaga by their mouthparts and manner of feeding. The mouthparts consist of three very slender stylets which form a tube. When at rest, they are retracted in a pocket which lies just behind the mouth. The antennae are usually five-segmented, rarely three-segmented. The thoracic segments are always very closely fused and lack wings. The claw is one-segmented, and on at least one pair of legs this claw is enlarged and can be folded into a process from the tibia to grasp a hair of the host. In some species two and in others all of the legs are thus modified. The ovipositor is very much reduced. It consists of but little more than two flaps which are able to close around a hair. Eyes are commonly lacking, and in those species in which they do occur, they are reduced to a pair of simple lenses or two light-receptive spots. All the species are quite small;
Anorexia nervosa the largest scarcely exceeds 0.2 in. (5 mm) in length and the smallest does not attain 0.04 in. (1 mm). Parasitic activity. The eggs are attached to the hairs of the host by a glue which issues from the base of the ovipositor and envelops a hair and one end of the egg. The young escape from the egg through a lid or operculum at the free end. In general, they are much like the adult and immediately begin feeding upon the blood of the host. Feeding is accomplished by thrusting forth the tube formed by the stylets, piercing the skin, and sucking blood by a pump in the throat of the insect. Very few species have been studied closely enough to determine their life-span, but in the human louse this is about a month. Disease transmission. This habit of sucking blood gives the Anoplura a special importance. They take up any disease-producing organisms in the host’s blood and may transfer these organisms to another individual. Thus, they transfer the organisms which cause two of the most important diseases of humans, epidemic typhus and relapsing fever. A third disease, trench fever, which gained great prominence during World War I, has also been of very great importance. The transmission of these diseases depends entirely upon certain habits of the insects, which are coincident with certain habits of their hosts. See RELAPSING FEVER. Host specificity. Indeed, these lice have become so closely adapted to life upon a particular kind of mammal, both physiologically and physically, that they can scarcely move about except when they are surrounded by its hair or body covering, and they can subsist only on its blood for any length of time. The transfer of a louse, from one host to another, can occur only when the hosts are in close bodily contact, as in the nest or at the time of mating. In the case of humans, the lice may become detached and enter the clothing or the bedding. In times of social disturbance, such as war or when people are crowded together as formerly occurred on ships, or in jails or slums, the opportunities for an exchange of these parasites are enormously increased. Under these conditions, in certain parts of the world louseborne diseases may become epidemic and produce tremendous mortality. The close restriction of each species of lice to its special host is sufficient to account for the fact that diseases are not transferred from one species of mammal to another. This has made it possible to prevent the epidemic spread of human diseases transmitted by lice by employing methods of cleanliness and destroying the lice. One epidemic disease of rodents that is carried by lice is known from India, and there are in all probability others. Human species. The degree to which lice are restricted to one kind of mammal or to a few closely related kinds may be indicated in the following discussion. There are two species of lice which occur as common ectoparasites upon humans. One of them is Pediculus humanus, the head and body louse which transmits the diseases mentioned above (Fig. 1). Apparently it sometimes occurs upon gibbons in zoos and has been recorded from certain
spiracle
antenna eye
head
coxa trochanter
thorax
femur tibia tarsus
abdomen
tarsal claw
(a)
(b)
Fig. 1. The louse Pediculus humanus, female, shown in (a) dorsal and (b) ventral views.
eye
tarsal claw
femur antenna
tibia tubercle
spiracle
(a)
(b)
Fig. 2. The louse Phthirius pubis, female, shown in (a) dorsal and (b) ventral views.
New World monkeys. Other quite distinct species of the genus Pediculus are known from the chimpanzee and from New World monkeys. The other species occurring on humans is Phthirius pubis, known as the crab louse, and it is restricted as far as is known to humans, although a second species has been described from the gorilla (Fig. 2). It transmits no known disease. The lice of the Old World monkeys are referred to another genus, Pedicinus. See INSECTA. Gordon F. Ferris; Dwight M. DeLong
Anorexia nervosa A psychiatric disorder in which a dramatic reduction in caloric intake consequent to excessive dieting leads to significant bodily, physiological,
713
714
Anorexia nervosa biochemical, emotional, psychological, and behavioral disturbances. Clinical presentation. Anorexia nervosa is typically an illness of adolescent females: 90% of all cases begin in girls who are between 12 and 20 years of age. Nevertheless, this disorder can also occur in males, in prepubertal girls, and in women well into their third decade. Moreover, if the illness becomes chronic, it can persist into mid-life and beyond. Estimates of the incidence of anorexia nervosa are limited but indicate that about 1–5% of adolescent girls living in industrialized cultures are affected, and that the rate is rising. Anorexia nervosa literally means “nervous loss of appetite” but appears to have little to do with such. Rather, there is usually a conscious decision made by a teen-age girl, most commonly around the ages of 14 or 18 years, to embark upon a diet. This decision often occurs in the context of some stressful life change situation. In addition, there are usually two personality characteristics: a morbid fear of losing control over one’s weight (the majority of girls who develop anorexia nervosa usually are slightly overweight prior to the illness’ onset) and a distorted body image (that is, the perception of oneself as larger than one is in reality). The amount of weight lost can vary considerably. The usual criterion for making a diagnosis of anorexia nervosa is a weight change of at least 25% from premorbid weight (to at least 15% below ideal weight in persons who were overweight at the onset). But this figure should be viewed as only a rule of thumb; some persons may be affected by the disorder and do not reach the 25% weight loss. In the most extreme cases, weight losses of up to 50% and more can occur. In addition to these core disturbances, there is an array of associated symptoms that characterize most persons with anorexia nervosa. Amenorrhea, or absence of menstruation, occurs in 95–100% of anorectic women; about one-third of the cases of amenorrhea actually precede the onset of substantial weight loss. Increased physical activity, such as excessive running, swimming, or calisthenics, is noticeable in a majority of cases; early in the course of the illness it is associated with a sense of well-being and is goal-directed, but later it becomes more obsessional and frenetic. Insomnia, whether secondary to the starvation or the emotional turmoil characteristic of the disorder, is a frequently reported symptom; in some cases, it can lead to an abuse of sedatives or tranquilizers. Use of emetics, cathartics, and diuretics is not uncommon in women who suffer from anorexia nervosa for at least several months. These practices are employed in an attempt to minimize the absorption of food, to rid the body of weight associated with body fluids, and to symbolically cleanse the self of “alien” substances. Difficulty in recognizing satiation is a common complaint, and if the illness becomes chronic, sensations of hunger may also become difficult to discriminate. Obsessional thinking and depression are not unusual in anorectic women; whether these symptoms precede the illness is not
always clear, but they become increasingly apparent with chronicity. The course and prognosis of anorexia nervosa is highly variable. While a high recovery rate, perhaps above 67%, is found in those persons whose illness begins acutely in their early teens and who quickly receive treatment, the outlook is considerably bleaker in those persons who develop the disorder later, who do not receive early treatment, and who develop bulimia. For all individuals with anorexia nervosa, the recovery rate is probably around 40%, with another 30% showing moderate improvement and the remaining 30% running a debilitated, chronic course. Included in the last group is the 5–10% of those who actually die from such complications of the disorder as cardiovascular collapse, infection, and suicide. Bulimia. Perhaps as many as 40–50% of anorectic patients whose illness persists for more than 1– 2 years will develop the additional eating disturbance known as bulimia. Persons whose anorexia starts relatively late, for example, in their twenties, may be particularly vulnerable to becoming bulimic. Literally meaning “oxhunger,” bulimia refers to binge eating or compulsive overeating wherein thousands of calories are consumed in a relatively brief period of time (for example, 2 h). The episodes usually occur in the evening in private and may be no more frequent than a few times per month. However, they can occur several times per day; the average rate is around several times per week. The binge characteristically involves carbohydrates—the very type of food that the anorectic so rigidly avoids the rest of the time—and will usually culminate in self-induced vomiting. There is usually enormous shame and selfloathing associated with bulimia; whereas dieting is equated with control and strength, binging has the opposite implication. Patients characteristically wake up the next morning with even greater resolve to diet, initiating a fast–feast cycle. The precise nature of the relationship of bulimia to anorexia nervosa remains unclear. Not all anorectics become bulimic and not all bulimics were anorectic. But dieting is common as a precursor to bulimia, and high premorbid weight and chronicity of weight loss seem to predispose the anorectic to developing bulimia. There is some indication that anorectic persons who develop bulimia are also more prone to depression (and often have a family history of depression), impulsive behavior (for example, shoplifting), and more outgoing but chaotic social relationships than those persons who remain “pure restricting” anorectics. A major concern about bulimia relates not to the overeating per se but to the subsequent vomiting. The chronic loss of gastric juices can lead to a depletion of the body’s potassium, which in turn can seriously affect cardiac function (to the point of death). In general, the onset of bulimia is believed to signify a poor prognosis. Psychological aspects. Although there is a broad range of symptoms, personality styles, precipitants, and outcomes that characterize anorexia nervosa, there is no simple explanation of its origins. In one
Anorexia nervosa widely accepted conception, the illness is viewed as a desperate struggle by the vulnerable female adolescent to establish a sense of identity separate from that of her domineering, overbearing, controlling, and intrusive mother. By virtue of the mother’s insensitivity to the cues of her child, such as hunger as an infant, the future anorectic has not received clear-cut confirmation of her internal experiences. This, coupled with the mother’s overintrusiveness and domination, leaves the developing child with serious deficits in her sense of self and in her ability to identify and discriminate feelings and bodily sensations. Her body image is disturbed, and she feels a potentially paralyzing sense of ineffectiveness. In the context of the typical stresses of puberty and adolescence, for example, bodily changes of portending adult responsibility and sexuality, separation from mother and home, the vicissitudes of relations with peers, social and sexual encounters with the opposite sex, and the demands of school or work—the preanorectic, who is already sensing high expectations from her achievement-oriented family, turns toward control of her body weight as the only area left in her life in which she can achieve fully independent and effective control. There is considerable emphasis on viewing the family system as the matrix in which the illness develops and for which the illness must have significance. Four characteristics of the families of anorectics are distinguishable: excessive mutual enmeshment, mutual overprotectiveness, rigidity in relation to internal problems and the external world, and an inability to achieve resolution of conflict. Anorexia is viewed as both a response to the lack of “living space” that the adolescent experiences and as a defense against the threats to the stability of the family system that the girl’s normal development implies. Development of anorexia is thus a function of both individual and family. It should be noted, however, that these formulations suffer from a common problem. They are based on assessments of anorectic patients—and their families—after the illness has become established. Thus, these theories cannot very well differentiate among predisposing, precipitating, and sustaining factors. Moreover, many of the problems and conflicts emphasized in anorectics—such as those concerning separation and independence, appearance, adult heterosexual relations—are no different from those of normally developing adolescents. And the more pathological features described in these individuals and their families—for example, feelings of emptiness, precarious sense of self, family enmeshment and rigidity—have been reported in other pathological but nonanorectic psychiatric conditions such as schizophrenia, borderline personality, and psychosomatic disorders. See SCHIZOPHRENIA; STRESS (PSYCHOLOGY). Biological aspects. There are numerous physical, physiological, and biochemical changes that reflect primarily, although not solely, the ravages of starvation. In addition to the general bodily emaciation they manifest, anorectic individuals show brit-
tle nails and thinning hair. Cold extremities, a slow pulse, a small heart, and a hypometabolic state, in compensation for the starvation, are also characteristic. In advanced cases, edema may occur. In anorectic women who also binge and vomit, tooth decay and enlargement of the salivary glands are common. A mild-to-moderate anemia and a diminished white blood cell count develop with progressive malnutrition. Diabetes insipidus, consequent to disturbed hypothalamic function, can also occur in advanced cases, producing large amounts of dilute urine and constant thirst. In addition, abnormalities in glucose tolerance and blood electrolytes have been observed. Particularly in women who vomit, the blood potassium level can be significantly low, possibly producing abnormalities on the electrocardiogram. A small percentage of anorectic subjects show vague abnormalities in the electroencephalogram. Interestingly, these abnormalities of electroencephalogram are not always clearly correlated with the extent of emaciation. See MALNUTRITION; PITUITARY GLAND DISORDERS. Because of the prominence of amenorrhea in its symptomatology, there has been a long-standing interest in the endocrinology of anorexia nervosa. A number of reliable hormonal abnormalities have been documented. Most prominent among these are diminished hypothalamic–pituitary–gonadal axis function and elevated hypothalamic–pituitary– adrenal axis function. While urinary and plasma gonadotropin levels are low and their circadian and monthly cycles are deficient, plasma and urinary levels of cortisol are high and usually will not be suppressed by dexamethasone (a synthetic adrenal steroid) administration. These abnormalities point to hypothalamic dysfunction, presumably consequent to weight loss and starvation. Possibly they also reflect functional changes secondary to affective disturbances (for example, depression), behavioral changes (for example, excessive exercise, altered sleep patterns), or some preexisting biological vulnerability. Reversal of these endocrine aberrations usually occurs with clinical improvement, although considerable time may be required for full normalization. See ENDOCRINE MECHANISMS. Treatment. As with virtually all psychiatric conditions for which the etiology is unknown and where no single empirically effective treatment exists, the therapeutic approaches to anorexia nervosa are diverse and reflect the different disciplines, training biases, and experiences of their proponents. Nevertheless, there are some general considerations in the management of this illness. First, there are probably mild cases that resolve without any intervention whatsoever. Second, despite this, anorexia nervosa is a potentially pernicious illness and, once established, can be strikingly resistant to all treatments; hence, early intervention is called for once the diagnosis is made. Third, the approach to treatment must be tailored to the stage of the illness. In the earliest cases and youngest children, the pediatrician’s involvement in a supportive, gently inquiring, and educational fashion may be sufficient.
715
716
Anorthite If no improvement is seen within a few weeks, the child may be referred to a psychiatrist, but the pediatrician should remain involved. Individual, insightoriented psychotherapy directed toward increasing confidence in identifying and accepting bodily feelings, understanding the sources of one’s low selfesteem and poor sense of self, and “trying out one’s wings” in the context of a reliable, supportive, respectful, and inquiring, but nonintrusive, relationship has generally been the essence of treatment for nonhospitalized individuals, particularly before chronicity has set in. Individual psychotherapy still remains a critical part of any approach to treatment, but the recovery rate can be increased, perhaps beyond 80%, by the inclusion of regular family therapy as part of the treatment approach. This implies the willingness of all the family members to participate and requires the individual to be young enough to still be involved with her family system. Hospitalization must be seriously considered whenever weight loss becomes substantial, exercising or food-related rituals begin to interfere with day-to-day functioning, or the individual begins to appear physically weak or emotionally hopeless. See PSYCHOTHERAPY. In the hospital, the first priority of treatment is directed toward correcting the biological abnormalities created by the extreme dieting (and, when present, the vomiting). If the individual appears unable or unwilling to resume adequate caloric intake, despite firm but supportive nursing and concomitant individual and family psychotherapy, more extreme measures may have to be instituted, including behavior modification or intravenous hyperalimentation. The behavior modification approach assumes that the food avoidance, whatever its original determinants, has become a habit that can best be treated by appropriate positive and negative reinforcement. Thus, a schedule is mapped out which progressively “punishes” weight loss with confinement to bed, loss of visitors, and so on, while “rewarding” weight gain with off-ward privileges, access to exercising, and such. While intravenous fluids can correct electrolyte imbalance, they cannot provide sufficient calories and nutrition to achieve weight gain. However, the availability now of total parenteral nutrition, that is, hyperalimentation with concentrated nutrients via a peripheral or a central large intravenous catheter, has provided a technique for reversing the severe emaciation caused by advanced anorexia nervosa. There is some evidence that, particularly in anorectics who are also characterized by depression and bulimia, antidepressant or possibly anticonvulsant medication may be helpful in damping down the binging and thereby gradually normalizing eating behavior in general. Neuroleptics (such as chlorpromazine) can also be helpful in the acute hospital management, particularly in the anxious or excessively exercising patient, but their beneficial effects are probably only indirect, that is, secondary to general tranquilization. See PSYCHOPHARMACOLOGY; TRANQUILIZER. Anorexia nervosa can become a chronic illness
for many of its victims. Improvement in the hospital does not mean that care need not be provided after discharge. Moreover, while a certain percentage of individuals do make a complete recovery, the real therapeutic challenge is in working with the chronic anorectic victim over many years in an attempt to minimize the morbidity imposed by the emaciation, medical complications, depression, and social isolation which often occur. See AFFECTIVE DISORDERS. Jack L. Katz Bibliography. H. Bruch, Eating Disorders: Obesity, Anorexia Nervosa, and the Person Within, 1979; H. Bruch, The Golden Cage: The Enigma of Anorexia Nervosa, 1978; P. E. Garfinkel and D. M. Garner, Anorexia Nervosa: A Multidimensional Perspective, 1982; S. Minuchin, B. Rosman, and L. Baker, Psychosomatic Families: Anorexia Nervosa in Context, 1978; J. A. Sours, Starving to Death in a Sea of Objects: The Anorexia Nervosa Syndrome, 1980.
Anorthite The calcium-rich plagioclase feldspar with composition Ab0An100 to Ab10An90 (Ab = NaAlSi3O8; An = CaAl2Si2O8), occurring in olivine-rich igneous rocks and rare volcanic ejecta (for example, at Mount Vesuvius, Italy, and Miyakejima, Japan). Hardness is 6 on Mohs scale, specific gravity 2.76, melting point 1550◦C (2822◦F). The crystal structure of Ab0An100 (triclinic space group P1¯ ) consists of an infinite three-dimensional array of corner-sharing [AlO4] and [SiO4] tetrahedra, alternately linked together in a framework of [Al2Si2O8]∞2− composition in which charge-balancing calcium (Ca2+) cations occupy four distinct, irregular cavities. The crystal structure of anorthite inverts to higher symmetry (I1¯ ) with each of the following: (1) heating above 242◦C (468◦F), (2) adding ∼10 mole % Ab component (substituting Na+ + Si4+ ↔ Ca2+ + Al3+), (3) increasing hydrostatic pressure above ∼2.6 gigapascals (26,000 atmospheres), and (4) effectively disordering the aluminum (Al) and silicon (Si) atoms at very high temperatures. Each inversion is evidenced by certain types of structurally out-of-step (antiphase) domains that can be imaged by transmission electron microscopy. Natural anorthite has no commercial uses, but the synthetic material [CaO · Al2O3 · 2SiO2 (known as CAS2)] is important in the ceramic industry and in certain composite materials with hightemperature applications. See CRYSTAL STRUCTURE; FELDSPAR; IGNEOUS ROCKS. Paul H. Ribbe
Anorthoclase The name usually given to alkali feldspars which have a chemical composition ranging from Or40Ab60 to Or10Ab90± up to approximately 20 mole % An (Or, Ab, An = KAlSi3O8, NaAlSi3O8, CaAl2Si2O8) and which deviate in one way or another from monoclinic symmetry tending toward triclinic symmetry. When found in nature, they usually do not consist
Anorthosite of a single phase but are composed of two or more kinds of K- and Na-rich domains mostly of submicroscopic size. In addition, they are frequently polysynthetically twinned after either or both of the albite and pericline laws. It appears that they originally grew as the monoclinic monalbite phase inverting and unmixing in the course of cooling during geological times. They are typically found in lavas or hightemperature rocks. See FELDSPAR; IGNEOUS ROCKS. Fritz H. Laves
Anorthosite A rock composed of 90 vol % or more of plagioclase feldspar. Strictly, the rock is composed entirely of crystals discernible with the eye, but some finely crystalline examples from the Moon have been called anorthosite or anorthositic breccia. Two principal types of anorthosite are based on field occurrence: layers in stratified complexes of igneous rock, and large massifs of rock up to 12,000 mi2 (30,000 km2) in area. Scientists have been fascinated with anorthosites because they are spectacular rocks (dark varieties are quarried and polished for ornamental use); valuable deposits of iron and titanium ore are associated with anorthosites; and the massif anorthosites appear to have been produced during a unique episode of ancient Earth history (about 1– 2 billion years ago). Definition, occurrence, and structure. Pure anorthosite has less than 10% of dark minerals— generally some combination of pyroxene, olivine, and oxides of iron and titanium; amphibole and biotite are rare, as are the light minerals apatite, zircon, scapolite, and calcite. Rocks with less than 90% but more than 78% of plagioclase are modified anorthosites (such as gabbroic anorthosite), and rocks with 78–65% of plagioclase are anorthositic (such as anorthositic gabbro). See GABBRO. The structure, texture, and mineralogy vary with type of occurrence. One type of occurrence is as layers (up to 10 ft or 3 m thick) interstratified with layers rich in pyroxene or olivine. The bulk composition of the layered rock samples containing anorthosite layers is gabbroic. The Stillwater igneous complex in Montana is an example where the lower part of the complex is rich in olivine, pyroxene, and chromite, and the upper part is rich in plagioclase. The second type of occurrence is the massif type. Commonly, the massifs are domical in shape and weakly layered. There appear to be two kinds of massif anorthosites: irregular massifs typified by the Lake St. John body in Quebec, and domical massifs similar to the Adirondack (New York) massif made classic through the studies of A. Buddington. The irregular massifs are older than the domical massifs and commonly contain olivine. The domical massifs are associated with silicic rocks of uncertain origin. See MASSIF. Possibly there is a third group of anorthosite occurrences: extremely ancient bodies of layered rock in which the layers of anorthosite contain calciumrich plagioclase and the adjacent layers are rich in
chromite and amphibole in addition to pyroxene. There are only a few examples of these apparently igneous complexes, in Greenland, southern Africa, and India. However, they appear to be terrestrial counterparts of lunar anorthosites. Recrystallization and deformation have pervasively affected these rocks, and it is difficult to discover their original attributes. Layered anorthosites. Zones of anorthosite in layered complexes are composed mostly of wellformed, moderately zoned crystals of plagioclase (labradorite to bytownite) up to several millimeters long. The tabular crystals of plagioclase commonly are aligned subparallel to the layering. Irregular crystals of pyroxene, olivine, and, rarely, magnetite are interstitial to the plagioclase crystals. These features are best explained by the hypothesis of accumulation of plagioclase crystals onto the floor or roof of the reservoir of melt, where gravity and perhaps motion of melt caused them to lie with their centers of gravity in the lowest position (like pennies spilled on the floor). After accumulation, the intergranular melt gradually crystallized as overgrowths of plagioclase and interstitial pyroxene. The separation of accumulated plagioclase crystals in anorthosite layers from pyroxene or olivine in pyroxenite and peridotite layers is a puzzle to scientists, for it is expected that plagioclase and pyroxene crystallize together in sub-equal proportions from basaltic magmas, not first one mineral, then the other and back again as suggested by the layering. Compositional zoning of plagioclase within layers and other structural and textural features suggest that plagioclase floated in certain intrusions. This is consistent with calculated densities of some postulated liquids, if the liquids are poor in water. Evolution of the magma body may be dynamically complex at the stage of plagioclase flotation and anorthosite formation, particularly if new batches of magma arriving from below are less dense than the old differentiated magma. Certain unusual rocks, including the platinum-rich Merensky reef horizon of the Bushveld intrusion located near the transition between mafic layers and anorthositic layers, may be related to such dynamical complexities. See MAGMA. Massif anorthosites. In massifs the plagioclase has an intermediate composition, mostly labradorite. The plagioclase crystals are poorly formed and practically unzoned and may be up to around 10 ft (3 m) long, but crystals about a centimeter long constitute the bulk of most massifs. Rare dark varieties have moderately well-formed plagioclase crystals (see illus.) filled with tiny specks of dark minerals. In many specimens the large labradorite crystals display a play of colors (blue and green are common) in a particular orientation, an attractive quality related to very slow cooling and sought after in ornamental stone. Lighter varieties have smaller crystals of plagioclase lacking inclusions of dark minerals, which instead occur as millimeter-sized interstitial grains. (Thus the various shades of anorthosite colors result from variations in the size or granularity of the pigmenting dark minerals.) In transitional varieties large dark crystals of plagioclase are set in a fine-grained
717
718
Anorthosite
Outcrop of massif anorthosite on the shore of Pipmuacan Reservoir, Quebec. Exceptionally well-preserved texture shows well-formed relict plagioclase crystals (now recrystallized) outlined by irregular black pyroxene grains. Two dark cores of twinned plagioclase are shown (circles).
sugarlike white matrix. Many of the plagioclase grains are bent and broken. The associated dark minerals are pyroxene, olivine, magnetite, and ilmenite. Deposits of almost pure magnetite occur and have been used for iron ore, although most deposits have too much titanium to be highly desired for iron ore. Other deposits are rich in olivine, spinel, and apatite as well as magnetite. The massifs are surrounded by gneissic rocks, and the original age relationships between anorthosite and gneiss are commonly uncertain. The textural and mineralogical features of massifs are indicative mainly of extensive recrystallization and deformation at high temperature. The dark varieties have textures similar to those of layered anorthosites and presumably have a similar igneous origin, but this cannot be determined with certainty. Many anorthosite massifs formed about 1.1– 1.7 billion years ago. Possibly they mark a unique episode in the history of the differentiation of the Earth. The restricted age of massif anorthosite still awaits a convincing explanation. Lunar anorthosites. By comparison with terrestrial occurrences, most lunar anorthosites are very fine grained, although one rock has crystals up to a centimeter long. Much of the fine grain size results from comminution by meteorite impact, and some of it probably results from rapid crystallization of impact melts. The unbrecciated lunar anorthosites have well-formed crystals of calcium-rich plagioclase (anorthite) intergrown with traces of olivine, spinel, pyroxene, ilmenite, metallic iron, and a few other exotic minerals. Associated rocks and minor minerals suggest that lunar anorthosites are of the layered type formed by accumulation of crystals of plagioclase from basaltic melts. Evidently, lunar plagioclases floated in iron-rich lunar basaltic melts, and most scientists agree that rocks rich in plagioclase and including anorthosite compose a major part of the skin of the Moon. See ANDESINE; LABRADORITE; MOON. Alfred T. Anderson, Jr. Role in growth of continental crust. Large occurrences of Archean anorthosite are formed principally
in back-arc basin or oceanic settings. Although intrusions of Archean anorthosite into continental crust are known, they seem limited in size, and lack of exposure further hampers evaluation of their contribution to continental growth. Based upon what is known of these rocks, it may be concluded that their contributions to Archean continents were vastly less than that of gray tonalitic gneiss and komatiite-basalt associates of the greenstone belts. See ARCHEAN. In contrast to the Archean, Proterozoic anorthosites were emplaced into continental crust in large massifs usually localized within belts such as the Grenville Province of Canada or the Rogaland-Ergusund region of southern Norway. In these belts anorthositic massifs may account for as much as 20% of the local crust. Most of these bodies were emplaced during the time period 1650–1000 million years ago (Ma) and appear to be related to the slow, anorogenic break-up of a mid-Proterozoic supercontinent. Within the Grenville Province the anorthosite massifs appear to span the 500-million-year interval from 1650 to 1150 Ma, and it is probable that the interval consists of anorthositic pulses at 1650, 1450, and 1150 Ma. These rocks account for a minimum surface area of close to 39,000 mi2 (100,000 km2), or about 13% of the crustal area of the Grenville Province. If this same concentration is assumed downward throughout the continental crust, a maximum value for anorthositic additions to continental crust may be computed. For a crust 24 mi (40 km) thick this results in 1 × 106 mi3 (4 × 106 km3) of anorthositic addition to the regional crust over 500 million years, or an average rate of crustal addition of 0.2 mi3/year (0.8 km3/year). This is about 150– 250 times less than the 0.3 mi3/year (1.2 km3/year) rate estimated for present-day arc accretion or the 0.2–0.4 mi3/year (1–2 km3/year) estimates for crustal accretion during the interval 1900–1700 Ma. The belt of anorthosite massifs in the Grenville Province is about 900 mi (1500 km) in length. This gives an accretion rate per mile (kilometer) of magmatic belt of about 2 mi3/106 years (5 km3/106 years) as compared with about 12 mi3/106 years (30 km3/106 years) per mile of magmatic belt for modern arcs. If the anorthosite massifs are continued downward for only 3 mi (5 km) into the crust, their rate of crustal addition becomes about 0.2 mi3/106 years (1 km3/106 years) per mile of magmatic belt. This is a small fraction of other crustal growth mechanisms, but it is far from insignificant. It may be concluded that within the mid-Proterozoic the emplacement of anorogenic, massif anorthosites constituted a small but significant means of continental crustal growth. See CONTINENTS, EVOLUTION OF; IGNEOUS ROCKS; METAMORPHISM; PROTEROZOIC. James McLelland Bibliography. L. D. Ashwal, Anorthosites, Minerals and Rocks, vol. 21, Springer-Verlag, 1993; R. B. Hargraves (ed.), Physics of Magmatic Processes, 1980; Y. Isachsen (ed.), Origin of Anorthosite and Related Rocks, N.Y. State Mus. Sci. Serv. Mem. 18, 1969; A. R. McBirney, Igneous Petrology, 2d ed., 1992; J. McLelland, Crustal growth associated
Anoxic zones with anorogenic, mid-Proterozoic anorthosite massifs of northeastern North America, Tectonophysics, 161:331–334, 1989; B. Mason and W. G. Melson, The Lunar Rocks, 1970; A. Reymer and G. Schubert, Phanerozoic addition sites to the continental crust and crustal growth, Tectonics, 3:63–77, 1984.
Anostraca An order of branchiopod crustaceans, known as fairy shrimps and, in some cases, brine shrimps (Artemia). These organisms range up to about 4 in. (100 mm) in length, but usually are much smaller. The trunk consists of 19 to 27 segments, of which usually the first 11, sometimes the first 17 or 19, bear limbs, plus a telson bearing flattened furcal rami (see illus.). The trunk limbs are foliaceous, and all are of the same basic type but differ a little among themselves; each limb is differentiated into a series of endites and is usually filtratory. By beating in metachronal rhythm they propel the animal forward and draw food particles toward it. Anostracans usually swim with their ventral surface uppermost. There is no carapace. The eyes are pedunculate. The antennae are large and modified as claspers in the male. The large mandibles are of the rolling, crushing type. The other mouthparts are very small. Most anostracans subsist on minute particles that they sieve from the water with their trunk limbs and pass forward to the mouthparts. A few large species become carnivores as they grow, seizing small crustaceans, including in some cases smaller species of fairy shrimps. Reproduction is usually bisexual, but some brine shrimps are parthenogenetic. Females pass eggs into a ventral brood pouch and then shed them. Eggs are drought-resistant. They hatch as nauplii which develop by a series of molts, gradually adding segments posteriorly.
antennulae
eye
antennae
thorax
legs
penes 2 mm
abdomen
telson cercopods Branchinecta paludosa, male, small specimen, lateral aspect.
Fairy shrimps frequent fresh water, usually temporary pools, in all regions of the world, but are also found in predator-free waters in Arctic and Antarctic regions. Brine shrimps also escape predation by living in hypersaline habitats such as salt lakes and salt pans to which they have become physiologically adapted. See BRANCHIOPODA. Geoffrey Fryer Bibliography. S. P. Parker (ed.), Synopsis and Classification of Living Organisms, 2 vols., 1982.
Anoxic zones Oxygen-depleted regions in marine environments. The dynamic steady state between oxygen supply and consumption determines the oxygen concentration. In regions where the rate of consumption equals or exceeds the rate of supply, seawater becomes devoid of oxygen and thus anoxic. In the open ocean, the only large regions that approach anoxic conditions are in the equatorial Pacific between 50 and 1000 m (165 and 3300 ft) depth and in the northern Arabian Sea and the Bay of Bengal in the Indian Ocean between 100 and 1000 m (330 and 3300 ft) depth. The Pacific oxygen-depleted region consists of vast tongues extending from Central America and Peru nearly to the middle of the ocean in some places. In parts of this zone, oxygen concentrations become very low, 15 micromoles per liter (atmospheric saturation is 200–300 µmol/L). Pore waters of marine sediments are sometimes anoxic a short distance below the sediment–water interface. The degree of oxygen consumption in sediment pore waters depends upon the amount of organic matter reaching the sediments and the rate of bioturbation (mixing of the surface sediment by benthic animals). In shallow regions (continental shelf and slope), pore waters are anoxic immediately below the sedimentwater interface; in relatively rapid sedimentation-rate areas of the deep sea, the pore waters are usually anoxic within a few centimeters of the interface; and in pore waters of slowly accumulating deep-sea sediments, oxygen may never become totally depleted. See MARINE SEDIMENTS. Restricted basins (areas where water becomes temporarily trapped) are often either permanently or intermittently anoxic. Here, the anoxia can be a consequence of natural conditions or anthropogenic causes such as eutrophication, or a combination of them. Classic examples are the Baltic Sea, Black Sea, Gulf of Mexico, Carioca Trench off the coast of Venezuela, and fiords on the Norwegian and British Columbia coasts. Lakes that receive large amounts of nutrients or organic matter (either from natural or human-produced sources) are often anoxic during the period of summer stratification. See BALTIC SEA; BLACK SEA; EUTROPHICATION; FIORD; GULF OF MEXICO. The chemistry of many elements dissolved in seawater (particularly the trace elements) is vastly changed by the presence or absence of oxygen. Since large areas of the ocean water mass are in contact with oxygen-depleted pore waters, the potential
719
Anoxic zones closed from the atmosphere, sulfate represents the largest oxidant reservoir in seawater (without sediments). In anoxic parts of the water column, such as in the Baltic Sea or Black Sea or as shown in Fig. 1a for a fiord in British Columbia, the concentrations of hydrogen sulfide, H2S, increase soon after oxygen and nitrate are depleted. In anoxic water or sediment O2, µmol/ liter 0
50
100
150
200
250
300 350
20 40 60
depth, m
exists for anoxic conditions to have a marked effect on the chemistry of the sea. See SEAWATER. Reduction-oxidation (redox) reactions. Organic matter is formed directly or indirectly through photosynthesis, which is driven by the energy of the Sun. Carbon dioxide, CO2, the stable form of combined carbon, is here converted to energy-rich organic carbon compounds. The required electrons for this reduction are provided by water (H2O) through the release of elemental oxygen. Organic carbon is a very unstable form of carbon and reacts with various oxidants, yielding back the energy required for its photosynthesis. The most common of these reactions is the oxidation of organic carbon. The electrons are transferred back from carbon to oxygen, yielding CO2 and H2O as final products. This reaction fuels animal life on land and in the sea. In the absence of oxygen, organic matter reacts with various other oxidants (that is, electron acceptors) to form CO2 and to gain energy. See PHOTOSYNTHESIS. The major redox reactions that occur in natural systems are given below. They are listed in the order of decreasing gain of free energy (2 and 3 are interchangeable). For simplicity, CH2O is used to represent organic carbon.
O2
80 100 120 140
ΣH2S
1. Oxygen reduction: 160 O2 + CH2 O −→ CO2 + H2 O
180
2. Nitrate reduction (denitrification): 4NO− + 5CH2 O + 4H+ −→ 5CO2 + 2N2 + 7H2 O 3
0
10
20
(a)
30 40 50 ΣH2 S, µmol/ liter
60
70
Fe2+, Mn2+, µmol/ liter
3. Manganese(IV) reduction: 0
2MnO2 + CH2 O + 4H+ −→ CO2 + 2Mn2+ + 3H2 O
1.0
2.0
3.0
4.0
5.0
6.0
120
4. Iron(III) reduction:
20
2Fe2 O3 + CH2 O + 8H+ −→ CO2 + 4Fe2+ + 5H2 O
5. Sulfate reduction:
2CH2 O −→ CH4 + CO2
The reaction with oxygen is the most energyyielding and continues until O2 is nearly entirely consumed. Nitrate, NO3−, and manganese, Mn(IV), reduction follow the consumption of oxygen. NO3− generally exists in low concentrations,