Advanced Engineering Mathematics, International Student Edition

International Student Edition

ADVANCED ENGINEERING MATHEMATICS

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ADVANCED ENGINEERING MATHEMATICS International Student Edition

PETER V. O’NEIL University of Alabama at Birmingham

Australia Canada Mexico Singapore Spain United Kingdom United States

Advanced Engineering Mathematics, International Student Edition by Peter V. O’Neil

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Contents

PART

1

Chapter 1

Ordinary Differential Equations 1 First-Order Differential Equations

3

1.1 Preliminary Concepts 3 1.1.1 General and Particular Solutions 3 1.1.2 Implicitly Defined Solutions 4 1.1.3 Integral Curves 5 1.1.4 The Initial Value Problem 6 1.1.5 Direction Fields 7 1.2 Separable Equations 11 1.2.1 Some Applications of Separable Differential Equations 14 1.3 Linear Differential Equations 22 1.4 Exact Differential Equations 26 1.5 Integrating Factors 33 1.5.1 Separable Equations and Integrating Factors 37 1.5.2 Linear Equations and Integrating Factors 37 1.6 Homogeneous, Bernoulli, and Riccati Equations 38 1.6.1 Homogeneous Differential Equations 38 1.6.2 The Bernoulli Equation 42 1.6.3 The Riccati Equation 43 1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories 1.7.1 Mechanics 46 1.7.2 Electrical Circuits 51 1.7.3 Orthogonal Trajectories 53 1.8 Existence and Uniqueness for Solutions of Initial Value Problems 58

Chapter 2

Second-Order Differential Equations

46

61

2.1 Preliminary Concepts 61 2.2 Theory of Solutions of y

+ p
xy
+ q
xy = f
x 62 2.2.1 The Homogeneous Equation y

+ p
xy
+ q
x = 0 64 2.2.2 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x 2.3 Reduction of Order 69 2.4 The Constant Coefficient Homogeneous Linear Equation 73 2.4.1 Case 1: A 2 − 4B > 0 73 2.4.2 Case 2: A 2 − 4B = 0 74

68

v

vi

Contents 2.4.3 Case 3: A 2 − 4B < 0 74 2.4.4 An Alternative General Solution in the Complex Root Case 75 2.5 Euler’s Equation 78 2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x 82 2.6.1 The Method of Variation of Parameters 82 2.6.2 The Method of Undetermined Coefficients 85 2.6.3 The Principle of Superposition 91 2.6.4 Higher-Order Differential Equations 91 2.7 Application of Second-Order Differential Equations to a Mechanical System 2.7.1 Unforced Motion 95 2.7.2 Forced Motion 98 2.7.3 Resonance 100 2.7.4 Beats 102 2.7.5 Analogy with an Electrical Circuit 103

Chapter 3

The Laplace Transform

107

3.1 Definition and Basic Properties 107 3.2 Solution of Initial Value Problems Using the Laplace Transform 3.3 Shifting Theorems and the Heaviside Function 120 3.3.1 The First Shifting Theorem 120 3.3.2 The Heaviside Function and Pulses 122 3.3.3 The Second Shifting Theorem 125 3.3.4 Analysis of Electrical Circuits 129 3.4 Convolution 134 3.5 Unit Impulses and the Dirac Delta Function 139 3.6 Laplace Transform Solution of Systems 144 3.7 Differential Equations with Polynomial Coefficients 150

Chapter 4

Series Solutions 4.1 4.2 4.3 4.4

Chapter 5

155

Power Series Solutions of Initial Value Problems 156 Power Series Solutions Using Recurrence Relations 161 Singular Points and the Method of Frobenius 166 Second Solutions and Logarithm Factors 173

Numerical Approximation of Solutions

181

5.1 Euler’s Method 182 5.1.1 A Problem in Radioactive Waste Disposal 5.2 One-Step Methods 190 5.2.1 The Second-Order Taylor Method 190 5.2.2 The Modified Euler Method 193 5.2.3 Runge-Kutta Methods 195 5.3 Multistep Methods 197 5.3.1 Case 1 r = 0 198 5.3.2 Case 2 r = 1 198 5.3.3 Case 3 r = 3 199 5.3.4 Case 4 r = 4 199

187

116

93

Contents

PART

2

Chapter 6

Vectors and Linear Algebra 201 Vectors and Vector Spaces 6.1 6.2 6.3 6.4 6.5

Chapter 7

203

The Algebra and Geometry of Vectors 203 The Dot Product 211 The Cross Product 217 The Vector Space R n 223 Linear Independence, Spanning Sets, and Dimension in R n

Matrices and Systems of Linear Equations

228

237

7.1 Matrices 238 7.1.1 Matrix Algebra 239 7.1.2 Matrix Notation for Systems of Linear Equations 242 7.1.3 Some Special Matrices 243 7.1.4 Another Rationale for the Definition of Matrix Multiplication 246 7.1.5 Random Walks in Crystals 247 7.2 Elementary Row Operations and Elementary Matrices 251 7.3 The Row Echelon Form of a Matrix 258 7.4 The Row and Column Spaces of a Matrix and Rank of a Matrix 266 7.5 Solution of Homogeneous Systems of Linear Equations 272 7.6 The Solution Space of AX = O 280 7.7 Nonhomogeneous Systems of Linear Equations 283 7.7.1 The Structure of Solutions of AX = B 284 7.7.2 Existence and Uniqueness of Solutions of AX = B 285 7.8 Matrix Inverses 293 7.8.1 A Method for Finding A−1 295

Chapter 8

Determinants 8.1 8.2 8.3 8.4 8.5 8.6 8.7 8.8 8.9

Chapter 9

299

Permutations 299 Definition of the Determinant 301 Properties of Determinants 303 Evaluation of Determinants by Elementary Row and Column Operations Cofactor Expansions 311 Determinants of Triangular Matrices 314 A Determinant Formula for a Matrix Inverse 315 Cramer’s Rule 318 The Matrix Tree Theorem 320

Eigenvalues, Diagonalization, and Special Matrices 9.1 Eigenvalues and Eigenvectors 324 9.1.1 Gerschgorin’s Theorem 328 9.2 Diagonalization of Matrices 330 9.3 Orthogonal and Symmetric Matrices

339

323

307

vii

viii

Contents 9.4 Quadratic Forms 347 9.5 Unitary, Hermitian, and Skew Hermitian Matrices

PART

3

Chapter 10

352

Systems of Differential Equations and Qualitative Methods 359 Systems of Linear Differential Equations

361

10.1 Theory of Systems of Linear First-Order Differential Equations 361 10.1.1 Theory of the Homogeneous System X
= AX 365 10.1.2 General Solution of the Nonhomogeneous System X
= AX + G 372 10.2 Solution of X
= AX when A is Constant 374 10.2.1 Solution of X
= AX when A has Complex Eigenvalues 377 10.2.2 Solution of X
= AX when A does not have n Linearly Independent Eigenvectors 379 10.2.3 Solution of X
= AX by Diagonalizing A 384 10.2.4 Exponential Matrix Solutions of X
= AX 386 10.3 Solution of X
= AX + G 394 10.3.1 Variation of Parameters 394 10.3.2 Solution of X
= AX + G by Diagonalizing A 398

Chapter 11

Qualitative Methods and Systems of Nonlinear Differential Equations 11.1 11.2 11.3 11.4 11.5 11.6 11.7

PART

4

Chapter 12

Nonlinear Systems and Existence of Solutions 403 The Phase Plane, Phase Portraits and Direction Fields Phase Portraits of Linear Systems 413 Critical Points and Stability 424 Almost Linear Systems 431 Lyapunov’s Stability Criteria 451 Limit Cycles and Periodic Solutions 461

406

Vector Analysis 473 Vector Differential Calculus

475

12.1 Vector Functions of One Variable 475 12.2 Velocity, Acceleration, Curvature and Torsion 481 12.2.1 Tangential and Normal Components of Acceleration 488 12.2.2 Curvature as a Function of t 491 12.2.3 The Frenet Formulas 492 12.3 Vector Fields and Streamlines 493 12.4 The Gradient Field and Directional Derivatives 499 12.4.1 Level Surfaces, Tangent Planes and Normal Lines 503 12.5 Divergence and Curl 510 12.5.1 A Physical Interpretation of Divergence 512 12.5.2 A Physical Interpretation of Curl 513

403

Contents

Chapter 13

Vector Integral Calculus

517

13.1 Line Integrals 517 13.1.1 Line Integral with Respect to Arc Length 525 13.2 Green’s Theorem 528 13.2.1 An Extension of Green’s Theorem 532 13.3 Independence of Path and Potential Theory in the Plane 536 13.3.1 A More Critical Look at Theorem 13.5 539 13.4 Surfaces in 3-Space and Surface Integrals 545 13.4.1 Normal Vector to a Surface 548 13.4.2 The Tangent Plane to a Surface 551 13.4.3 Smooth and Piecewise Smooth Surfaces 552 13.4.4 Surface Integrals 553 13.5 Applications of Surface Integrals 557 13.5.1 Surface Area 557 13.5.2 Mass and Center of Mass of a Shell 557 13.5.3 Flux of a Vector Field Across a Surface 560 13.6 Preparation for the Integral Theorems of Gauss and Stokes 562 13.7 The Divergence Theorem of Gauss 564 13.7.1 Archimedes’s Principle 567 13.7.2 The Heat Equation 568 13.7.3 The Divergence Theorem as a Conservation of Mass Principle 13.8 The Integral Theorem of Stokes 572 13.8.1 An Interpretation of Curl 576 13.8.2 Potential Theory in 3-Space 576

PART

5

Chapter 14

570

Fourier Analysis, Orthogonal Expansions, and Wavelets 581 Fourier Series

583

14.1 Why Fourier Series? 583 14.2 The Fourier Series of a Function 586 14.2.1 Even and Odd Functions 589 14.3 Convergence of Fourier Series 593 14.3.1 Convergence at the End Points 599 14.3.2 A Second Convergence Theorem 601 14.3.3 Partial Sums of Fourier Series 604 14.3.4 The Gibbs Phenomenon 606 14.4 Fourier Cosine and Sine Series 609 14.4.1 The Fourier Cosine Series of a Function 610 14.4.2 The Fourier Sine Series of a Function 612 14.5 Integration and Differentiation of Fourier Series 614 14.6 The Phase Angle Form of a Fourier Series 623 14.7 Complex Fourier Series and the Frequency Spectrum 630 14.7.1 Review of Complex Numbers 630 14.7.2 Complex Fourier Series 631

ix

x

Contents

Chapter 15

The Fourier Integral and Fourier Transforms 15.1 15.2 15.3 15.4

15.5 15.6 15.7

15.8

15.9

Chapter 16

637

The Fourier Integral 637 Fourier Cosine and Sine Integrals 640 The Complex Fourier Integral and the Fourier Transform 642 Additional Properties and Applications of the Fourier Transform 652 15.4.1 The Fourier Transform of a Derivative 652 15.4.2 Frequency Differentiation 655 15.4.3 The Fourier Transform of an Integral 656 15.4.4 Convolution 657 15.4.5 Filtering and the Dirac Delta Function 660 15.4.6 The Windowed Fourier Transform 661 15.4.7 The Shannon Sampling Theorem 665 15.4.8 Lowpass and Bandpass Filters 667 The Fourier Cosine and Sine Transforms 670 The Finite Fourier Cosine and Sine Transforms 673 The Discrete Fourier Transform 675 15.7.1 Linearity and Periodicity 678 15.7.2 The Inverse N -Point DFT 678 15.7.3 DFT Approximation of Fourier Coefficients 679 Sampled Fourier Series 681 15.8.1 Approximation of a Fourier Transform by an N -Point DFT 685 15.8.2 Filtering 689 The Fast Fourier Transform 694 15.9.1 Use of the FFT in Analyzing Power Spectral Densities of Signals 15.9.2 Filtering Noise From a Signal 696 15.9.3 Analysis of the Tides in Morro Bay 697

Special Functions, Orthogonal Expansions, and Wavelets

695

701

16.1 Legendre Polynomials 701 16.1.1 A Generating Function for the Legendre Polynomials 704 16.1.2 A Recurrence Relation for the Legendre Polynomials 706 16.1.3 Orthogonality of the Legendre Polynomials 708 16.1.4 Fourier–Legendre Series 709 16.1.5 Computation of Fourier–Legendre Coefficients 711 16.1.6 Zeros of the Legendre Polynomials 713 16.1.7 Derivative and Integral Formulas for Pn
x 715 16.2 Bessel Functions 719 16.2.1 The Gamma Function 719 16.2.2 Bessel Functions of the First Kind and Solutions of Bessel’s Equation 16.2.3 Bessel Functions of the Second Kind 722 16.2.4 Modified Bessel Functions 725 16.2.5 Some Applications of Bessel Functions 727 16.2.6 A Generating Function for Jn
x 732 16.2.7 An Integral Formula for Jn
x 733 16.2.8 A Recurrence Relation for Jv
x 735 16.2.9 Zeros of Jv
x 737

721

Contents 16.2.10 Fourier–Bessel Expansions 739 16.2.11 Fourier–Bessel Coefficients 741 16.3 Sturm–Liouville Theory and Eigenfunction Expansions 745 16.3.1 The Sturm–Liouville Problem 745 16.3.2 The Sturm–Liouville Theorem 752 16.3.3 Eigenfunction Expansions 755 16.3.4 Approximation in the Mean and Bessel’s Inequality 759 16.3.5 Convergence in the Mean and Parseval’s Theorem 762 16.3.6 Completeness of the Eigenfunctions 763 16.4 Wavelets 765 16.4.1 The Idea Behind Wavelets 765 16.4.2 The Haar Wavelets 767 16.4.3 A Wavelet Expansion 774 16.4.4 Multiresolution Analysis with Haar Wavelets 774 16.4.5 General Construction of Wavelets and Multiresolution Analysis 16.4.6 Shannon Wavelets 776

PART

6

Chapter 17

xi

775

Partial Differential Equations 779 The Wave Equation

781

17.1 The Wave Equation and Initial and Boundary Conditions 781 17.2 Fourier Series Solutions of the Wave Equation 786 17.2.1 Vibrating String with Zero Initial Velocity 786 17.2.2 Vibrating String with Given Initial Velocity and Zero Initial Displacement 791 17.2.3 Vibrating String with Initial Displacement and Velocity 793 17.2.4 Verification of Solutions 794 17.2.5 Transformation of Boundary Value Problems Involving the Wave Equation 796 17.2.6 Effects of Initial Conditions and Constants on the Motion 798 17.2.7 Numerical Solution of the Wave Equation 801 17.3 Wave Motion Along Infinite and Semi-Infinite Strings 808 17.3.1 Wave Motion Along an Infinite String 808 17.3.2 Wave Motion Along a Semi-Infinite String 813 17.3.3 Fourier Transform Solution of Problems on Unbounded Domains 815 17.4 Characteristics and d’Alembert’s Solution 822 17.4.1 A Nonhomogeneous Wave Equation 825 17.4.2 Forward and Backward Waves 828 17.5 Normal Modes of Vibration of a Circular Elastic Membrane 831 17.6 Vibrations of a Circular Elastic Membrane, Revisited 834 17.7 Vibrations of a Rectangular Membrane 837

Chapter 18

The Heat Equation

841

18.1 The Heat Equation and Initial and Boundary Conditions 18.2 Fourier Series Solutions of the Heat Equation 844

841

xii

Contents 18.2.1 Ends of the Bar Kept at Temperature Zero 844 18.2.2 Temperature in a Bar with Insulated Ends 847 18.2.3 Temperature Distribution in a Bar with Radiating End 848 18.2.4 Transformations of Boundary Value Problems Involving the Heat Equation 851 18.2.5 A Nonhomogeneous Heat Equation 854 18.2.6 Effects of Boundary Conditions and Constants on Heat Conduction 857 18.2.7 Numerical Approximation of Solutions 859 18.3 Heat Conduction in Infinite Media 865 18.3.1 Heat Conduction in an Infinite Bar 865 18.3.2 Heat Conduction in a Semi-Infinite Bar 868 18.3.3 Integral Transform Methods for the Heat Equation in an Infinite Medium 869 18.4 Heat Conduction in an Infinite Cylinder 873 18.5 Heat Conduction in a Rectangular Plate 877

Chapter 19

The Potential Equation

879

19.1 19.2 19.3 19.4 19.5

Harmonic Functions and the Dirichlet Problem 879 Dirichlet Problem for a Rectangle 881 Dirichlet Problem for a Disk 883 Poisson’s Integral Formula for the Disk 886 Dirichlet Problems in Unbounded Regions 888 19.5.1 Dirichlet Problem for the Upper Half Plane 889 19.5.2 Dirichlet Problem for the Right Quarter Plane 891 19.5.3 An Electrostatic Potential Problem 893 19.6 A Dirichlet Problem for a Cube 896 19.7 The Steady-State Heat Equation for a Solid Sphere 898 19.8 The Neumann Problem 902 19.8.1 A Neumann Problem for a Rectangle 904 19.8.2 A Neumann Problem for a Disk 906 19.8.3 A Neumann Problem for the Upper Half Plane 908

PART

7

Chapter 20

Complex Analysis 911 Geometry and Arithmetic of Complex Numbers

913

20.1 Complex Numbers 913 20.1.1 The Complex Plane 914 20.1.2 Magnitude and Conjugate 915 20.1.3 Complex Division 916 20.1.4 Inequalities 917 20.1.5 Argument and Polar Form of a Complex Number 918 20.1.6 Ordering 920 20.2 Loci and Sets of Points in the Complex Plane 921 20.2.1 Distance 922 20.2.2 Circles and Disks 922 20.2.3 The Equation z − a = z − b 923 20.2.4 Other Loci 925 20.2.5 Interior Points, Boundary Points, and Open and Closed Sets

925

Contents 20.2.6 20.2.7 20.2.8 20.2.9

Chapter 21

Limit Points 929 Complex Sequences 931 Subsequences 934 Compactness and the Bolzano-Weierstrass Theorem

Complex Functions

935

939

21.1 Limits, Continuity, and Derivatives 939 21.1.1 Limits 939 21.1.2 Continuity 941 21.1.3 The Derivative of a Complex Function 943 21.1.4 The Cauchy–Riemann Equations 945 21.2 Power Series 950 21.2.1 Series of Complex Numbers 951 21.2.2 Power Series 952 21.3 The Exponential and Trigonometric Functions 957 21.4 The Complex Logarithm 966 21.5 Powers 969 21.5.1 Integer Powers 969 21.5.2 z 1/n for Positive Integer n 969 21.5.3 Rational Powers 971 21.5.4 Powers z w 972

Chapter 22

Complex Integration

975

22.1 Curves in the Plane 975 22.2 The Integral of a Complex Function 980 22.2.1 The Complex Integral in Terms of Real Integrals 983 22.2.2 Properties of Complex Integrals 985 22.2.3 Integrals of Series of Functions 988 22.3 Cauchy’s Theorem 990 22.3.1 Proof of Cauchy’s Theorem for a Special Case 993 22.4 Consequences of Cauchy’s Theorem 994 22.4.1 Independence of Path 994 22.4.2 The Deformation Theorem 995 22.4.3 Cauchy’s Integral Formula 997 22.4.4 Cauchy’s Integral Formula for Higher Derivatives 1000 22.4.5 Bounds on Derivatives and Liouville’s Theorem 1001 22.4.6 An Extended Deformation Theorem 1002

Chapter 23

Series Representations of Functions

1007

23.1 Power Series Representations 1007 23.1.1 Isolated Zeros and the Identity Theorem 1012 23.1.2 The Maximum Modulus Theorem 1016 23.2 The Laurent Expansion 1019

Chapter 24

Singularities and the Residue Theorem

1023

24.1 Singularities 1023 24.2 The Residue Theorem 1030 24.3 Some Applications of the Residue Theorem

1037

xiii

xiv

Contents 24.3.1 The Argument Principle 1037 24.3.2 An Inversion for the Laplace Transform 24.3.3 Evaluation of Real Integrals 1040

Chapter 25

Conformal Mappings 25.1 25.2 25.3 25.4 25.5

PART

8

Chapter 26

1039

1055

Functions as Mappings 1055 Conformal Mappings 1062 25.2.1 Linear Fractional Transformations 1064 Construction of Conformal Mappings Between Domains 1072 25.3.1 Schwarz–Christoffel Transformation 1077 Harmonic Functions and the Dirichlet Problem 1080 25.4.1 Solution of Dirichlet Problems by Conformal Mapping 1083 Complex Function Models of Plane Fluid Flow 1087

Probability and Statistics 1097 Counting and Probability

1099

26.1 26.2 26.3

The Multiplication Principle 1099 Permutations 1102 Choosing r Objects from n Objects 1104 26.3.1 r Objects from n Objects, with Order 1104 26.3.2 r Objects from n Objects, without Order 1106 26.3.3 Tree Diagrams 1107 26.4 Events and Sample Spaces 1112 26.5 The Probability of an Event 1116 26.6 Complementary Events 1121 26.7 Conditional Probability 1122 26.8 Independent Events 1126 26.8.1 The Product Rule 1128 26.9 Tree Diagrams in Computing Probabilities 1130 26.10 Bayes’ Theorem 1134 26.11 Expected Value 1139

Chapter 27

Statistics 27.1

27.2 27.3

27.4

1143

Measures of Center and Variation 1143 27.1.1 Measures of Center 1143 27.1.2 Measures of Variation 1146 Random Variables and Probability Distributions 1150 The Binomial and Poisson Distributions 1154 27.3.1 The Binomial Distribution 1154 27.3.2 The Poisson Distribution 1157 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve 27.4.1 The Standard Bell Curve 1174 27.4.2 The 68, 95, 99.7 Rule 1176

1159

Contents 27.5 27.6 27.7 27.8

Sampling Distributions and the Central Limit Theorem 1178 Confidence Intervals and Estimating Population Proportion 1185 Estimating Population Mean and the Student t Distribution 1190 Correlation and Regression 1194

Answers and Solutions to Selected Problems Index

I1

A1

xv

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Preface

This Sixth Edition of Advanced Engineering Mathematics maintains the primary goal of previous editions—to engage much of the post-calculus mathematics needed and used by scientists, engineers, and applied mathematicians, in a setting that is helpful to both students and faculty. The format used throughout begins with the correct developments of concepts such as Fourier series and integrals, conformal mappings, and special functions. These ideas are then brought to bear on applications and models of important phenomena, such as wave and heat propagation and filtering of signals. This edition differs from the previous one primarily in the inclusion of statistics and numerical methods. The statistics part treats random variables, normally distributed data, bell curves, the binomial, Poisson, and student t-distributions, the central limit theorem, confidence intervals, correlation, and regression. This is preceded by prerequisite topics from probability and techniques of enumeration. The numerical methods are applied to initial value problems in ordinary differential equations, including a proposal for radioactive waste disposal, and to boundary value problems involving the heat and wave equations. Finally, in order to include these topics without lengthening the book, some items from the fifth edition have been moved to a website, located at http://engineering.thomsonlearning.com. I hope that this provides convenient accessibility. Material selected for this move includes some biographies and historical notes, predator/prey and competing species models, the theory underlying the efficiency of the FFT, and some selected examples and problems. The chart on the following page offers a complete organizational overview.

Acknowledgments This book is the result of a team effort involving much more than an author. Among those to whom I owe a debt of appreciation are Chris Carson, Joanne Woods, Hilda Gowans and Kamilah Reid-Burrell of Thomson Engineering, and Rose Kernan and the professionals at RPK Editorial Services, Inc. I also want to thank Dr. Thomas O’Neil of the California Polytechnic State University for material he contributed, and Rich Jones, who had the vision for the first edition of this book many years ago. Finally, I want to acknowledge the reviewers, whose suggestions for improvements and clarifications are much appreciated: Preliminary Review

Panagiotis Dimitrakopoulos, University of Maryland Mohamed M. Hafez, University of California, Davis Jennifer Hopwood, University of Western Australia Nun Kwan Yip, Purdue University xvii

xviii Organizational Overview Ordinary Differential Equations

Laplace Transforms

Series Solutions

Systems of Ordinary Differential Equations

Special Functions

Vectors, Matrices, Determinants

Statistical Analysis Systems of Algebraic Equations

Eigenfunction Expansions, Completeness Haar Wavelets

Vector Analysis Qualitative Methods, Stability, Analysis of Critical Points Probability

Partial Differential Equations Statistics

Mathematical Models

Fourier Analysis Complex Analysis Fourier Series, Integrals

Fourier Transforms

Discrete Fourier Transform

P r eface

xix

Draft Review

Sabri Abou-Ward, University of Toronto Craig Hildebrand, California State University – Fresno Seiichi Nomura, University of Texas, Arlington David L. Russell, Virginia Polytechnic Institute and State University Y.Q. Sheng, McMaster University

Peter V. O’neil University of Alabama at Birmingham

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PA RT

1

CHAPTER 1 First-Order Differential Equations

CHAPTER 2 Second-Order Differential Equations

CHAPTER 3

Ordinary Differential Equations

The Laplace Transform

CHAPTER 4 Series Solutions

CHAPTER 5 Numerical Approximation of Solutions

A differential equation is an equation that contains one or more derivatives. For example, y


x + y
x = 4 sin
3x and d4 w −
w
t2 = e−t dt4 are differential equations. These are ordinary differential equations because they involve only total derivatives, rather than partial derivatives. Differential equations are interesting and important because they express relationships involving rates of change. Such relationships form the basis for developing ideas and studying phenomena in the sciences, engineering, economics, and increasingly in other areas, such as the business world and the stock market. We will see examples of applications as we learn more about differential equations. 1

The order of a differential equation is the order of its highest derivative. The first example given above is of second order, while the second is of fourth order. The equation xy
− y2 = ex is of first order. A solution of a differential equation is any function that satisfies it. A solution may be defined on the entire real line, or on only part of it, often an interval. For example, y = sin
2x is a solution of y

+ 4y = 0 because, by direct differentiation, y

+ 4y = −4 sin
2x + 4 sin
2x = 0 This solution is defined for all x (that is, on the whole real line). By contrast, y = x ln
x − x is a solution of y
=

y + 1 x

but this solution is defined only for x > 0. Indeed, the coefficient 1/x of y in this equation means that x = 0 is disallowed from the start. We now begin a systematic development of ordinary differential equations, starting with the first order case.

2

CHAPTER

1

PRELIMINARY CONCEPTS SEPARABLE EQUATIONS HOMOGENEOUS, BERNOULLI, AND RICCATI EQUATIONS APPLICATIONS TO MECHANICS, ELECTRICAL CIRCUITS, AND ORTHOGONAL TRAJECTORIES EXI

First-Order Differential Equations

1.1

Preliminary Concepts Before developing techniques for solving various kinds of differential equations, we will develop some terminology and geometric insight.

1.1.1

General and Particular Solutions

A first-order differential equation is any equation involving a first derivative, but no higher derivative. In its most general form, it has the appearance F
x y y
 = 0

(1.1)

in which y
x is the function of interest and x is the independent variable. Examples are y
− y2 − ey = 0 y
− 2 = 0 and y
− cos
x = 0 Note that y
must be present for an equation to qualify as a first-order differential equation, but x and/or y need not occur explicitly. A solution of equation (1.1) on an interval I is a function  that satisfies the equation for all x in I. That is, F
x 
x 

x = 0

for all x in I.

For example, 
x = 2 + ke−x 3

4

CHAPTER 1

First-Order Differential Equations

is a solution of y
+ y = 2 for all real x, and for any number k. Here I can be chosen as the entire real line. And 
x = x ln
x + cx is a solution of y
=

y +1 x

for all x > 0, and for any number c. In both of these examples, the solution contained an arbitrary constant. This is a symbol independent of x and y that can be assigned any numerical value. Such a solution is called the general solution of the differential equation. Thus 
x = 2 + ke−x is the general solution of y
+ y = 2. Each choice of the constant in the general solution yields a particular solution. For example, f
x = 2 + e−x  and

g
x = 2 − e−x

√ h
x = 2 − 53e−x

are√all particular solutions of y
+ y = 2, obtained by choosing, respectively, k = 1, −1 and − 53 in the general solution.

1.1.2

Implicitly Defined Solutions

Sometimes we can write a solution explicitly giving y as a function of x. For example, y = ke−x is the general solution of y
= −y as can be verified by substitution. This general solution is explicit, with y isolated on one side of an equation, and a function of x on the other. By contrast, consider y
= −

2xy3 + 2  3x2 y2 + 8e4y

We claim that the general solution is the function y
x implicitly defined by the equation x2 y3 + 2x + 2e4y = k

(1.2)

in which k can be any number. To verify this, implicitly differentiate equation (1.2) with respect to x, remembering that y is a function of x. We obtain 2xy3 + 3x2 y2 y
+ 2 + 8e4y y
= 0 and solving for y
yields the differential equation. In this example we are unable to solve equation (1.2) explicitly for y as a function of x, isolating y on one side. Equation (1.2), implicitly defining the general solution, was obtained by a technique we will develop shortly, but this technique cannot guarantee an explicit solution.

1.1 Preliminary Concepts

1.1.3

5

Integral Curves

A graph of a solution of a first-order differential equation is called an integral curve of the equation. If we know the general solution, we obtain an infinite family of integral curves, one for each choice of the arbitrary constant.

EXAMPLE 1.1

We have seen that the general solution of y
+ y = 2 is y = 2 + ke−x for all x. The integral curves of y
+ y = 2 are graphs of y = 2 + ke−x for different choices of k. Some of these are shown in Figure 1.1. y 30

k6

20 k3 10 k  0 ( y  2)
2 k 
3

k 
6


1

0

1

2

3


10


20

FIGURE 1.1 Integral curves of y
+ y = 2 for k = 0 3 −3 6, and −6.

EXAMPLE 1.2

It is routine to verify that the general solution of y
+

y = ex x

is y=

1
xex − ex + c x

4

5

6

x

6

CHAPTER 1

First-Order Differential Equations

for x = 0. Graphs of some of these integral curves, obtained by making choices for c, are shown in Figure 1.2. y 40 30

c  20

20 10

c5 c0

0.5
10
20

1.5

1.0

2.0

2.5

x 3.0

c 
6 c 
10

FIGURE 1.2

−10.

Integral curves of y
+ x1 y = ex for c = 0 5 20 −6, and

We will see shortly how these general solutions are obtained. For the moment, we simply want to illustrate integral curves. Although in simple cases integral curves can be sketched by hand, generally we need computer assistance. Computer packages such as MAPLE, MATHEMATICA and MATLAB are widely available. Here is an example in which the need for computing assistance is clear.

EXAMPLE 1.3

The differential equation y
+ xy = 2 has general solution y
x = e−x

2 /2




x

2e

0

2 /2

d + ke−x /2  2

Figure 1.3 shows computer-generated integral curves corresponding to k = 0, 4, 13, −7, −15 and −11.

1.1.4

The Initial Value Problem

The general solution of a first-order differential equation F
x y y
 = 0 contains an arbitrary constant, hence there is an infinite family of integral curves, one for each choice of the constant. If we specify that a solution is to pass through a particular point
x0  y0 , then we must find that particular integral curve (or curves) passing through this point. This is called an initial value problem. Thus, a first order initial value problem has the form F
x y y
 = 0

y
x0  = y0 

in which x0 and y0 are given numbers. The condition y
x0  = y0 is called an initial condition.

1.1 Preliminary Concepts

7

y k  13 10 k4

5
2

k0

0

k 
7


4

2

4

x


5 k 
11


10

k 
15
15 FIGURE 1.3 Integral curves of y
+ xy = 2 for k = 0 4 13 −7 −15, and

−11.

EXAMPLE 1.4

Consider the initial value problem y
+ y = 2

y
1 = −5




From Example 1.1, the general solution of y + y = 2 is y = 2 + ke−x  Graphs of this equation are the integral curves. We want the one passing through
1 −5. Solve for k so that y
1 = 2 + ke−1 = −5 obtaining k = −7e The solution of this initial value problem is y = 2 − 7ee−x = 2 − 7e−
x−1  As a check, y
1 = 2 − 7 = −5 The effect of the initial condition in this example was to pick out one special integral curve as the solution sought. This suggests that an initial value problem may be expected to have a unique solution. We will see later that this is the case, under mild conditions on the coefficients in the differential equation.

1.1.5

Direction Fields

Imagine a curve, as in Figure 1.4. If we choose some points on the curve and, at each point, draw a segment of the tangent to the curve there, then these segments give a rough outline of the shape of the curve. This simple observation is the key to a powerful device for envisioning integral curves of a differential equation.

8

CHAPTER 1

First-Order Differential Equations y

x FIGURE 1.4 Short tangent segments suggest the shape of the curve.

The general first-order differential equation has the form F
x y y
 = 0 Suppose we can solve for y
and write the differential equation as y
= f
x y Here f is a known function. Suppose f
x y is defined for all points
x y in some region R of the plane. The slope of the integral curve through a given point
x0  y0  of R is y

x0 , which equals f
x0  y0 . If we compute f
x y at selected points in R, and draw a small line segment having slope f
x y at each
x y, we obtain a collection of segments which trace out the shapes of the integral curves. This enables us to obtain important insight into the behavior of the solutions (such as where solutions are increasing or decreasing, limits they might have at various points, or behavior as x increases). A drawing of the plane, with short line segments of slope f
x y drawn at selected points
x y, is called a direction field of the differential equation y
= f
x y. The name derives from the fact that at each point the line segment gives the direction of the integral curve through that point. The line segments are called lineal elements.

EXAMPLE 1.5

Consider the equation y
= y2  Here f
x y = y2 , so the slope of the integral curve through
x y is y2 . Select some points and, through each, draw a short line segment having slope y2 . A computer generated direction field is shown in Figure 1.5(a). The lineal elements form a profile of some integral curves and give us some insight into the behavior of solutions, at least in this part of the plane. Figure 1.5(b) reproduces this direction field, with graphs of the integral curves through
0 1,
0 2,
0 3,
0 −1,
0 −2 and
0 −3. By a method we will develop, the general solution of y
= y2 is y=−

1  x+k

so the integral curves form a family of hyperbolas, as suggested by the curves sketched in Figure 1.5(b).

1.1 Preliminary Concepts

9

y 4

2


4


2

2

0

4

x


2


4 FIGURE 1.5(a) A direction field for y
= y2 .

y 4

2

2
4


2

4

0


2


4 FIGURE 1.5(b) Direction field for y
= y2 and integral curves through
0 1,
0 2,
0 3
0 −1,
0 −2, and
0 −3.

x

10

CHAPTER 1

First-Order Differential Equations

EXAMPLE 1.6

Figure 1.6 shows a direction field for y
= sin
xy together with the integral curves through
0 1,
0 2,
0 3,
0 −1,
0 −2 and
0 −3. In this case, we cannot write a simple expression for the general solution, and the direction field provides information about the behavior of solutions that is not otherwise readily apparent. y 4

2


4


2

0

2

4

x


2


4 Direction field for y
= sin
xy and integral curves through
0 1,
0 2,
0 3,
0 −1,
0 −2, and
0 −3.

FIGURE 1.6

With this as background, we will begin a program of identifying special classes of firstorder differential equations for which there are techniques for writing the general solution. This will occupy the next five sections.

SECTION 1.1

PROBLEMS

In each of Problems 1 through 6, determine whether the given function is a solution of the differential equation. √ 1. 2yy
= 1 
x = x − 1 for x > 1 2. y
+ y = 0 
x = Ce−x C −e 2y + e for x > 0 
x = 2x 2x √ 2xy C 4. y
= for x = ± 2 
x = 2 2 − x2 x −2

3. y
= −

x

x

5. xy
= x − y 
x =

x2 − 3 for x = 0 2x

6. y
+ y = 1 
x = 1 + Ce−x In each of Problems 7 through 11, verify by implicit differentiation that the given equation implicitly defines a solution of the differential equation. 7. y2 + xy − 2x2 − 3x − 2y = C y − 4x − 2 +
x + 2y − 2y
= 0

1.2 Separable Equations 8. xy3 − y = C y3 +
3xy2 − 1y
= 0

17. y
= x + y y
2 = 2


18. y
= x − xy y
0 = −1

9. y − 4x + e = C 8x − ye −
2y + xe y = 0 2

2

xy

xy

xy

19. y
= xy y
0 = 2

10. 8 lnx − 2y + 4 − 2x + 6y = C; y
=

20. y
= x − y + 1 y
0 = 1

x − 2y 3x − 6y + 4

2x3 + 2xy2 − y x + 2 y
= 0 11. tan
y/x + x = C 2 2 x +y x + y2 −1

2

In each of Problems 12 through 16, solve the initial value problem and graph the solution. Hint: Each of these differential equations can be solved by direct integration. Use the initial condition to solve for the constant of integration. 12. y
= 2x y
2 = 1


In each of Problems 21 through 26, generate a direction field and some integral curves for the differential equation. Also draw the integral curve representing the solution of the initial value problem. These problems should be done by a software package. 21. y
= sin
y y
1 = /2 22. y
= x cos
2x − y y
1 = 0 23. y
= y sin
x − 3x2  y
0 = 1

−x

13. y = e  y
0 = 2

24. y
= ex − y y
−2 = 1




14. y = 2x + 2 y
−1 = 1

25. y
− y cos
x = 1 − x2  y
2 = 2

15. y
= 4 cos
xsin
x y
/2 = 0

26. y
= 2y + 3 y
0 = 1




16. y = 8x + cos
2x y
0 = −3 In each of Problems 17 through 20 draw some lineal elements of the differential equation for −4 ≤ x ≤ 4, −4 ≤ y ≤ 4. Use the resulting direction field to sketch a graph of the solution of the initial value problem. (These problems can be done by hand.)

1.2

11

27. Show that, for the differential equation y
+ p
xy = q
x, the lineal elements on any vertical line x = x0 , with p
x0  = 0, all pass through the single point
 , where  = x0 +

1 p
x0 

and

=

q
x0   p
x0 

Separable Equations DEFINITION 1.1

Separable Differential Equation

A differential equation is called separable if it can be written y
= A
xB
y

In this event, we can separate the variables and write, in differential form, 1 dy = A
x dx B
y wherever B
y = 0. We attempt to integrate this equation, writing



1 dy = A
x dx B
y

This yields an equation in x, y, and a constant of integration. This equation implicitly defines the general solution y
x. It may or may not be possible to solve explicitly for y
x.

12

CHAPTER 1

First-Order Differential Equations

EXAMPLE 1.7

y
= y2 e−x is separable. Write dy = y2 e−x dx as 1 dx = e−x dx y2 for y = 0. Integrate this equation to obtain 1 − = −e−x + k y an equation that implicitly defines the general solution. In this example we can explicitly solve for y, obtaining the general solution y=

1  −k

e−x

Now recall that we required that y = 0 in order to separate the variables by dividing by y2 . In fact, the zero function y
x = 0 is a solution of y
= y2 ex , although it cannot be obtained from the general solution by any choice of k. For this reason, y
x = 0 is called a singular solution of this equation. Figure 1.7 shows graphs of particular solutions obtained by choosing k as 0, 3, −3, 6 and −6. y 3

2 k0 1 k 
3 k 
6
2 k6


1 k3

1

2

x


1

Integral curves of y
= y2 e−x for k = 0 3 −3 6, and −6. FIGURE 1.7

Whenever we use separation of variables, we must be alert to solutions potentially lost through conditions imposed by the algebra used to make the separation.

1.2 Separable Equations

13

EXAMPLE 1.8

x2 y
= 1 + y is separable, and we can write 1 1 dy = 2 dx 1+y x The algebra of separation has required that x = 0 and y = −1, even though we can put x = 0 and y = −1 into the differential equation to obtain the correct equation 0 = 0. Now integrate the separated equation to obtain 1 ln1 + y = − + k x This implicitly defines the general solution. In this case, we can solve for y
x explicitly. Begin by taking the exponential of both sides to obtain 1 + y = ek e−1/x = Ae−1/x  in which we have written A = ek . Since k could be any number, A can be any positive number. Then 1 + y = ±Ae−1/x = Be−1/x  in which B = ±A can be any nonzero number. The general solution is y = −1 + Be−1/x  in which B is any nonzero number. Now revisit the assumption that x = 0 and y = −1. In the general solution, we actually obtain y = −1 if we allow B = 0. Further, the constant function y
x = −1 does satisfy x2 y
= 1 + y. Thus, by allowing B to be any number, including 0, the general solution y
x = −1 + Be−1/x contains all the solutions we have found. In this example, y = −1 is a solution, but not a singular solution, since it occurs as a special case of the general solution. Figure 1.8 shows graphs of solutions corresponding to B = −8 −5 0 4 and 7. y

B7

4

B4

2 0
2

1

2


4
6

3

4 5 B0

x

B 
5 B 
8


8 FIGURE 1.8 Integral curves of x2 y
= 1 + y for

B = 0 4 7 −5, and −8.

We often solve an initial value problem by finding the general solution of the differential equation, then solving for the appropriate choice of the constant.

14

CHAPTER 1

First-Order Differential Equations

EXAMPLE 1.9

Solve the initial value problem y
= y2 e−x  y
1 = 4 We know from Example 1.7 that the general solution of y
= y2 e−x is y
x =

1  −k

e−x

Now we need to choose k so that y
1 =

1 = 4 e−1 − k

from which we get 1 k = e−1 −  4 The solution of the initial value problem is y
x =

1 e−x

+ 41 − e−1



EXAMPLE 1.10

The general solution of y
= y


x − 12 y+3

is implicitly defined by 1 y + 3 lny =
x − 13 + k 3

(1.3)

To obtain the solution satisfying y
3 = −1, put x = 3 and y = −1 into equation (1.3) to obtain 1 −1 =
23 + k 3 hence k=−

11  3

The solution of this initial value problem is implicitly defined by 11 1 y + 3 lny =
x − 13 −  3 3

1.2.1

Some Applications of Separable Differential Equations

Separable equations arise in many contexts, of which we will discuss three.

1.2 Separable Equations

15

EXAMPLE 1.11

(The Mathematical Policewoman) A murder victim is discovered, and a lieutenant from the forensic science laboratory is summoned to estimate the time of death. The body is located in a room that is kept at a constant 68 degrees Fahrenheit. For some time after the death, the body will radiate heat into the cooler room, causing the body’s temperature to decrease. Assuming (for want of better information) that the victim’s temperature was a “normal” 986 at the time of death, the lieutenant will try to estimate this time by observing the body’s current temperature and calculating how long it would have had to lose heat to reach this point. According to Newton’s law of cooling, the body will radiate heat energy into the room at a rate proportional to the difference in temperature between the body and the room. If T
t is the body temperature at time t, then for some constant of proportionality k, T

t = k T
t − 68  The lieutenant recognizes this as a separable differential equation and writes 1 dT = k dt T − 68 Upon integrating, she gets lnT − 68 = kt + C Taking exponentials, she gets T − 68 = ekt+C = Aekt  in which A = eC . Then T − 68 = ±Aekt = Bekt  Then T
t = 68 + Bekt  Now the constants k and B must be determined, and this requires information. The lieutenant arrived at 9:40 p.m. and immediately measured the body temperature, obtaining 944 degrees. Letting 9:40 be time zero for convenience, this means that T
0 = 944 = 68 + B and so B = 264. Thus far, T
t = 68 + 264ekt  To determine k, the lieutenant makes another measurement. At 11:00 she finds that the body temperature is 892 degrees. Since 11:00 is 80 minutes past 9:40, this means that T
80 = 892 = 68 + 264e80k  Then e80k = so

212  264 

80k = ln

212 264



16

CHAPTER 1

First-Order Differential Equations

and

  212 1 ln  k= 80 264

The lieutenant now has the temperature function: T
t = 68 + 264eln
212/264t/80  In order to find when last time when the body was 986 (presumably the time of death), solve for the time in T
t = 986 = 68 + 264eln
212/264t/80  To do this, the lieutenant writes 306 = eln
212/264t/80 264 and takes the logarithm of both sides to obtain     212 306 t ln ln =  264 80 264 Therefore the time of death, according to this mathematical model, was 80 ln
306/264  ln
212/264

t=

which is approximately −538 minutes. Death occurred approximately 538 minutes before (because of the negative sign) the first measurement at 9:40, which was chosen as time zero. This puts the murder at about 8:46 p.m.

EXAMPLE 1.12

(Radioactive Decay and Carbon Dating) In radioactive decay, mass is converted to energy by radiation. It has been observed that the rate of change of the mass of a radioactive substance is proportional to the mass itself. This means that, if m
t is the mass at time t, then for some constant of proportionality k that depends on the substance, dm = km dt This is a separable differential equation. Write it as 1 dm = k dt m and integrate to obtain lnm = kt + c Since mass is positive, m = m and ln
m = kt + c Then m
t = ekt+c = Aekt  in which A can be any positive number.

1.2 Separable Equations

17

Determination of A and k for a given element requires two measurements. Suppose at some time, designated as time zero, there are M grams present. This is called the initial mass. Then m
0 = A = M so m
t = Mekt  If at some later time T we find that there are MT grams, then m
T = MT = MekT  Then

 ln

hence

MT M

 = kT

  MT 1  k = ln T M

This gives us k and determines the mass at any time: m
t = Meln
MT /Mt/T  We obtain a more convenient formula for the mass if we choose the time of the second measurement more carefully. Suppose we make the second measurement at that time T = H at which exactly half of the mass has radiated away. At this time, half of the mass remains, so MT = M/2 and MT /M = 1/2. Now the expression for the mass becomes m
t = Meln
1/2t/H  or m
t = Me− ln
2t/H 

(1.4)

This number H is called the half-life of the element. Although we took it to be the time needed for half of the original amount M to decay, in fact, between any times t1 and t1 + H, exactly half of the mass of the element present at t1 will radiate away. To see this, write m
t1 + H = Me− ln
2
t1 +H/H = Me− ln
2t1 /H e− ln
2H/H = e− ln
2 m
t1  = 21 m
t1  Equation (1.4) is the basis for an important technique used to estimate the ages of certain ancient artifacts. The earth’s upper atmosphere is constantly bombarded by high-energy cosmic rays, producing large numbers of neutrons, which collide with nitrogen in the air, changing some of it into radioactive carbon-14, or 14 C. This element has a half-life of about 5,730 years. Over the relatively recent period of the history of this planet in which life has evolved, the fraction of 14 C in the atmosphere, compared to regular carbon, has been essentially constant. This means that living matter (plant or animal) has injested 14 C at about the same rate over a long historical period, and objects living, say, two million years ago would have had the same ratio of carbon-14 to carbon in their bodies as objects alive today. When an organism dies, it ceases its intake of 14 C, which then begins to decay. By measuring the ratio of 14 C to carbon in an artifact, we can estimate the amount of the decay, and hence the time it took, giving an

18

CHAPTER 1

First-Order Differential Equations

estimate of the time the organism was alive. This process of estimating the age of an artifact is called carbon dating. Of course, in reality the ratio of 14 C in the atmosphere has only been approximately constant, and in addition a sample may have been contaminated by exposure to other living organisms, or even to the air, so carbon dating is a sensitive process that can lead to controversial results. Nevertheless, when applied rigorously and combined with other tests and information, it has proved a valuable tool in historical and archeological studies. To apply equation (1.4) to carbon dating, use H = 5730 and compute ln
2 ln
2 = ≈ 0000120968 H 5730 in which ≈ means “approximately equal” (not all decimal places are listed). Equation (1.4) becomes m
t = Me−0000120968t  Now suppose we have an artifact, say a piece of fossilized wood, and measurements show that the ratio of 14 C to carbon in the sample is 37 percent of the current ratio. If we say that the wood died at time 0, then we want to compute the time T it would take for one gram of the radioactive carbon to decay this amount. Thus, solve for T in 037 = e−0000120968T  We find that T =−

ln
037 ≈ 8,219 0000120968

years. This is a little less than one and one-half half-lives, a reasonable estimate if nearly the 14 C has decayed.

2 3

of

EXAMPLE 1.13

(Torricelli’s Law) Suppose we want to estimate how long it will take for a container to empty by discharging fluid through a drain hole. This is a simple enough problem for, say, a soda can, but not quite so easy for a large oil storage tank or chemical facility. We need two principles from physics. The first is that the rate of discharge of a fluid flowing through an opening at the bottom of a container is given by dV = −kAv dt in which V
t is the volume of fluid in the container at time t, v
t is the discharge velocity of fluid through the opening, A is the cross sectional area of the opening (assumed constant), and k is a constant determined by the viscosity of the fluid, the shape of the opening, and the fact that the cross-sectional area of fluid pouring out of the opening is slightly less than that of the opening itself. In practice, k must be determined for the particular fluid, container, and opening, and is a number between 0 and 1. We also need Torricelli’s law, which states that v
t is equal to the velocity of a free-falling particle released from a height equal to the depth of the fluid at time t. (Free-falling means that the particle is influenced by gravity only). Now the work done by gravity in moving the particle from its initial point by a distance h
t is mgh
t, and this must equal the change in the kinetic energy,
21 mv2 . Therefore,  v
t = 2gh
t

1.2 Separable Equations 18
h

r

19

18

h

FIGURE 1.9

Putting these two equations together yields  dV = −kA 2gh
t dt

(1.5)

We will apply equation (1.5) to a specific case to illustrate its use. Suppose we have a hemispherical tank of water, as in Figure 1.9. The tank has radius 18 feet, and water drains through a circular hole of radius 3 inches at the bottom. How long will it take the tank to empty? Equation (1.5) contains two unknown functions, V
t and h
t, so one must be eliminated. Let r
t be the radius of the surface of the fluid at time t and consider an interval of time from t0 to t1 = t0 + t. The volume V of water draining from the tank in this time equals the volume of a disk of thickness h (the change in depth) and radius r
t∗ , for some t∗ between t0 and t1 . Therefore V =  r
t∗ 2 h so h V =  r
t∗ 2  t t In the limit as t → 0, dh dV = r 2  dt dt Putting this into equation (1.5) yields r 2

 dh = −kA 2gh dt

Now V has been eliminated, but at the cost of introducing r
t. However, from Figure 1.9, r 2 = 182 −
18 − h2 = 36h − h2 so   dh  = −kA 2gh  36h − h2 dt This is a separable differential equation, which we write as 

 36h − h2 dh = −kA 2g dt 1/2 h

Take g to be 32 feet per second per second. The radius of the circular opening is 3 inches, or 1 feet, so its area is A = /16 square feet. For water, and an opening of this shape and size, 4 the experiment gives k = 08. The last equation becomes  √   1 1/2 3/2 dh = −
08 36h − h 64 dt 16

CHAPTER 1

20

First-Order Differential Equations

or



 36h1/2 − h3/2 dh = −04 dt

A routine integration yields 2 2 24h3/2 − h5/2 = − t + c 5 5 or 60h3/2 − h5/2 = −t + k Now h
0 = 18, so √

60
183/2 −
185/2 = k

Thus k = 2268 2 and h
t is implicitly determined by the equation √ 60h3/2 − h5/2 = 2268 2 − t √ The tank is empty when h = 0, and this occurs when t = 2268 2 seconds, or about 53 minutes, 28 seconds. The last three examples contain an important message. Differential equations can be used to solve a variety of problems, but a problem usually does not present itself as a differential equation. Normally we have some event or process, and we must use whatever information we have about it to derive a differential equation and initial conditions. This process is called mathematical modeling. The model consists of the differential equation and other relevant information, such as initial conditions. We look for a function satisfying the differential equation and the other information, in the hope of being able to predict future behavior, or perhaps better understand the process being considered.

SECTION 1.2

PROBLEMS

In each of Problems 1 through 10, determine if the differential equation is separable. If it is, find the general solution (perhaps implicitly defined). If it is not separable, do not attempt a solution at this time.

In each of Problems 11 through 15, solve the initial value problem. 11. xy2 y
= y + 1 y
3e2  = 2 12. y
= 3x2
y + 2 y
2 = 8

1. 3y
= 4x/y2

13. ln
yx y
= 3x2 y y
2 = e3

2. y + xy
= 0

14. 2yy
= ex−y  y
4 = −2

3. cos
yy
= sin
x + y

15. yy
= 2x sec
3y y
2/3 = /3

4. e

x+y


y = 3x




5. xy + y = y 6. y
=

2


x + 12 − 2y 2y

7. x sin
yy
= cos
y 8.

x
2y2 + 1 y = y x+1

9. y + y
= ex − sin
y 10. cos
x + y + sin
x − yy
= cos
2x

2

16. An object having a temperature of 90 degrees Fahrenheit is placed into an environment kept at 60 degrees. Ten minutes later the object has cooled to 88 degrees. What will be the temperature of the object after it has been in this environment for 20 minutes? How long will it take for the object to cool to 65 degrees? 17. A thermometer is carried outside a house whose ambient temperature is 70 degrees Fahrenheit. After five minutes the thermometer reads 60 degrees, and fifteen minutes after this, 504 degrees. What is the outside temperature (which is assumed to be constant)?

1.2 Separable Equations 18. Assume that the population of bacteria in a petri dish changes at a rate proportional to the population at that time. This means that, if P
t is the population at time t, then dP = kP dt for some constant k. A particular culture has a population density of 100,000 bacteria per square inch. A culture that covered an area of 1 square inch at 10:00 a.m. on Tuesday was found to have grown to cover 3 square inches by noon the following Thursday. How many bacteria will be present at 3:00 p.m. the following Sunday? How many will be present on Monday at 4:00 p.m.? When will the world be overrun by these bacteria, assuming that they can live anywhere on the earth’s surface? (Here you need to look up the land area of the earth.) 19. Assume that a sphere of ice melts at a rate proportional to its surface area, retaining a spherical shape. Interpret melting as a reduction of volume with respect to time. Determine an expression for the volume of the ice at any time t. 20. A radioactive element has a half-life of ln
2 weeks. If e3 tons are present at a given time, how much will be left 3 weeks later? 21. The half-life of uranium-238 is approximately 45 · 109 years. How much of a 10-kilogram block of U-238 will be present 1 billion years from now? 22. Given that 12 grams of a radioactive element decays to 91 grams in 4 minutes, what is the half-life of this element? 23. Evaluate




e−t

2 −9/t 2

dt

0

Hint: Let I
x =






e−t

2 −
x/t2

dt

0

Calculate I

x by differentiating under the integral sign, then let u = x/t. Show that I

x = −2I
x and solve for I
x. Evaluate by using the stan 2 the constant √ dard result that 0 e−t dt = /2. Finally, evaluate I
3. 24.  Derive the fact used in Example 1.13 that v
t = 2gh
t. Hint: Consider a free-falling particle having height h
t at time t. The work done by gravity in moving the particle from its starting point to a given point is mgh
t, and this must equal the change in the kinetic energy, which is
1/2mv2 .

21

25. Calculate the time required to empty the hemispherical tank of Example 1.13 if the tank is positioned with its flat side down. 26. (Draining a Hot Tub) Consider a cylindrical hot tub with a 5-foot radius and height of 4 feet, placed on one of its circular ends. Water is draining from the tub through a circular hole 58 inches in diameter located in the base of the tub. (a) Assume a value k = 06 to determine the rate at which the depth of the water is changing. Here it is useful to write dh dV dV/dt dh = =  dt dV dt dV/dh (b) Calculate the time T required to drain the hot tub if it is initially full. Hint: One way to do this is to write T=




0

H

dt dh dh

(c) Determine how much longer it takes to drain the lower half than the upper half of the tub. Hint: Use the integral suggested in (b), with different limits for the two halves. 27. (Draining a Cone) A tank shaped like a right circular cone, with its vertex down, is 9 feet high and has a diameter of 8 feet. It is initially full of water. (a) Determine the time required to drain the tank through a circular hole of diameter 2 inches at the vertex. Take k = 06. (b) Determine the time it takes to drain the tank if it is inverted and the drain hole is of the same size and shape as in (a), but now located in the new base. 28. (Drain Hole at Unknown Depth) Determine the rate of change of the depth of water in the tank of Problem 27 (vertex at the bottom) if the drain hole is located in the side of the cone 2 feet above the bottom of the tank. What is the rate of change in the depth of the water when the drain hole is located in the bottom of the tank? Is it possible to determine the location of the drain hole if we are told the rate of change of the depth and the depth of the water in the tank? Can this be done without knowing the size of the drain opening? 29. Suppose the conical tank of Problem 27, vertex at the bottom, is initially empty and water is added at the constant rate of /10 cubic feet per second. Does the tank ever overflow? 30. (Draining a Sphere) Determine the time it takes to completely drain a spherical tank of radius 18 feet if it is initially full of water and the water drains through a circular hole of radius 3 inches located in the bottom of the tank. Use k = 08.

22

1.3

CHAPTER 1

First-Order Differential Equations

Linear Differential Equations DEFINITION 1.2

Linear Differential Equation

A first-order differential equation is linear if it has the form y

x + p
xy = q
x Assume that p and q are continuous on an interval I (possibly the whole real line). Because of the special form of the linear equation, we can obtain the general solution on I by a clever  observation. Multiply the differential equation by e p
x dx to get 

e

p
x dx


y
x + p
xe



p
x dx

y = q
xe



p
x dx

The left side of this equation is the derivative of the product y
xe





p
x dx

, enabling us to write

  d  y
xe p
x dx = q
xe p
x dx  dx Now integrate to obtain 

y
xe

p
x dx

Finally, solve for y
x: y
x = e−



p
x dx


  q
xe p
x dx dx + C

=






q
xe

p
x dx



dx + Ce−



p
x dx



(1.6)



The function e p
x dx is called an integrating factor for the differential equation, because multiplication of the differential equation by this factor results in an equation that can be integrated to obtain the general solution. We do not recommend memorizing equation (1.6). Instead, recognize  the form of the linear equation and understand the technique of solving it by multiplying by e p
x dx . EXAMPLE 1.14

The equation y
+ y = sin
x is linear. Here p
x = 1 and q
x = sin
x, both continuous for all x. An integrating factor is 

e

dx



or ex . Multiply the differential equation by ex to get y
ex + yex = ex sin
x or
yex 
= ex sin
x Integrate to get


1 yex = ex sin
x dx = ex sin
x − cos
x + C 2 The general solution is y
x =

1

sin
x − cos
x + Ce−x  2

1.3 Linear Differential Equations

23

EXAMPLE 1.15

Solve the initial value problem y y
= 3x2 −  x

y
1 = 5

First recognize that the differential equation can be written in linear form: 1 y
+ y = 3x2  x An integrating factor is e to get




1/x dx

= eln
x = x, for x > 0. Multiply the differential equation by x xy
+ y = 3x3 

or
xy
= 3x3  Integrate to get 3 xy = x4 + C 4 Then 3 C y
x = x3 + 4 x for x > 0. For the initial condition, we need y
1 = 5 =

3 +C 4

so C = 17/4 and the solution of the initial value problem is 3 17 y
x = x3 + 4 4x for x > 0. Depending on p and q, it may not be possible to evaluate all of the integrals in the general solution 1.6 in closed form (as a finite algebraic combination of elementary functions). This occurs with y
+ xy = 2 whose general solution is y
x = 2e−x

2 /2




ex

2 /2

dx + Ce−x /2  2

 2 We cannot write ex /2 dx in elementary terms. However, we could still use a software package to generate a direction field and integral curves, as is done in Figure 1.10. This provides some idea of the behavior of solutions, at least within the range of the diagram.

24

CHAPTER 1

First-Order Differential Equations y 4

2


4


2

x 0

2

4


2


4

Integral curves of y
+ xy = 2 passing through
0 2,
0 4,
0 −2, and
0 −5. FIGURE 1.10

Linear differential equations arise in many contexts. Example 1.11, involving estimation of time of death, involved a separable differential equation which is also linear and could have been solved using an integrating factor.

EXAMPLE 1.16

(A Mixing Problem) Sometimes we want to know how much of a given substance is present in a container in which various substances are being added, mixed, and removed. Such problems are called mixing problems, and they are frequently encountered in the chemical industry and in manufacturing processes. As an example, suppose a tank contains 200 gallons of brine (salt mixed with water), in which 100 pounds of salt are dissolved. A mixture consisting of 18 pound of salt per gallon is flowing into the tank at a rate of 3 gallons per minute, and the mixture is continuously stirred. Meanwhile, brine is allowed to empty out of the tank at the same rate of 3 gallons per minute (Figure 1.11). How much salt is in the tank at any time?

1

8 lb/gal; 3 gal/min

3 gal/min

FIGURE 1.11

Before constructing a mathematical model, notice that the initial ratio of salt to brine in the tank is 100 pounds per 200 gallons, or 21 pound per gallon. Since the mixture pumped in has a constant ratio of 18 pound per gallon, we expect the brine mixture to dilute toward the incoming ratio, with a “terminal” amount of salt in the tank of 18 pound per gallon, times 200 gallons. This leads to the expectation that in the long term (as t → ) the amount of salt in the tank should approach 25 pounds.

1.3 Linear Differential Equations

25

Now let Q
t be the amount of salt in the tank at time t. The rate of change of Q
t with time must equal the rate at which salt is pumped in, minus the rate at which it is pumped out. Thus dQ = (rate in) − (rate out) dt       gallons Q
t pounds gallons 1 pounds 3 − 3 = 8 gallon minute 200 gallon minute =

3 3 − Q
t 8 200

This is the linear equation Q

t + An integrating factor is e obtain




3/200 dt

3 3 Q=  200 8

= e3t/200 . Multiply the differential equation by this factor to

Q
e3t/200 +

3 3t/200 3 e Q = e3t/200  200 8

or  3t/200 
3 3t/200 Qe = e  8 Then Qe3t/200 =

3 200 3t/200 e + C 8 3

so Q
t = 25 + Ce−3t/200  Now Q
0 = 100 = 25 + C so C = 75 and Q
t = 25 + 75e−3t/200  As we expected, as t increases, the amount of salt approaches the limiting value of 25 pounds. From the derivation of the differential equation for Q
t, it is apparent that this limiting value depends on the rate at which salt is poured into the tank, but not on the initial amount of salt in the tank. The term 25 in the solution is called the steady-state part of the solution because it is independent of time, and the term 75e−3t/200 is the transient part. As t increases, the transient part exerts less influence on the amount of salt in the tank, and in the limit the solution approaches its steady-state part.

CHAPTER 1

26

First-Order Differential Equations

SECTION 1.3

PROBLEMS

In each of Problems 1 through 8, find the general solution. Not all integrals can be done in closed form. 3 1. y
− y = 2x2 x 2. y
− y = sinh
x 3. y
+ 2y = x 4. sin
2xy
+ 2y sin2
x = 2 sin
x 5. y
− 2y = −8x2 6.
x2 − x − 2y
+ 3xy = x2 − 4x + 4 7. y
+ y =

x−1 x2

8. y
+ sec
xy = cos
x In each of Problems 9 through 14, solve the initial value problem. 9. y
+

1 y = 3x y
3 = 4 x−2

10. y
+ 3y = 5e2x − 6 y
0 = 2 11. y
+

2 y = 3 y
0 = 5 x+1

16. A 500-gallon tank initially contains 50 gallons of brine solution in which 28 pounds of salt have been dissolved. Beginning at time zero, brine containing 2 pounds of salt per gallon is added at the rate of 3 gallons per minute, and the mixture is poured out of the tank at the rate of 2 gallons per minute. How much salt is in the tank when it contains 100 gallons of brine? Hint: The amount of brine in the tank at time t is 50 + t. 17. Two tanks are cascaded as in Figure 1.12. Tank 1 initially contains 20 pounds of salt dissolved in 100 gallons of brine, while tank 2 contains 150 gallons of brine in which 90 pounds of salt are dissolved. At time zero a brine solution containing 21 pound of salt per gallon is added to tank 1 at the rate of 5 gallons per minute. Tank 1 has an output that discharges brine into tank 2 at the rate of 5 gallons per minute, and tank 2 also has an output of 5 gallons per minute. Determine the amount of salt in each tank at any time t. Also determine when the concentration of salt in tank 2 is a minimum and how much salt is in the tank at that time. Hint: Solve for the amount of salt in tank 1 at time t first and then use this solution to determine the amount in tank 2.

5 gal/min

12.
x2 − 2xy
+
x2 − 5x + 4y =
x4 − 2x3 e−x  y
3 = 18e−3

1 2

lb/gal Tank 1

13. y
− y = 2e4x  y
0 = −3 14. y
+

5y = 3x3 + x y
−1 = 4 9x

Tank 2 5 gal/min

15. Find all functions with the property that the y-intercept of the tangent to the graph at
x y is 2x2 .

1.4

5 gal/min

FIGURE 1.12

Mixing between tanks in Problem 17.

Exact Differential Equations We continue the theme of identifying certain kinds of first-order differential equations for which there is a method leading to a solution. We can write any first order equation y
= f
x y in the form M
x y + N
x yy
= 0 For example, put M
x y = −f
x y and N
x y = 1. An interesting thing happens if there is a function  such that

 = M
x y and

x

 = N
x y

y

(1.7)

1.4 Exact Differential Equations

27

In this event, the differential equation becomes

  dy + = 0

x y dx which, by the chain rule, is the same as d 
x y
x = 0 dx But this means that 
x y
x = C with C constant. If we now read this argument from the last line back to the first, the conclusion is that the equation 
x y = C implicitly defines a function y
x that is the general solution of the differential equation. Thus, finding a function that satisfies equation (1.7) is equivalent to solving the differential equation. Before taking this further, consider an example.

EXAMPLE 1.17

The differential equation y
= −

2xy3 + 2 3x2 y2 + 8e4y

is neither separable nor linear. Write it in the form   M + Ny
= 2xy3 + 2 + 3x2 y2 + 8e4y y
= 0

(1.8)

with M
x y = 2xy3 + 2

and

N
x y = 3x2 y2 + 8e4y 

Equation (1.8) can in turn be written M dx + N dy =
2xy3 + 2 dx +
3x2 y2 + 8e4y  dy = 0 Now let 
x y = x2 y3 + 2x + 2e4y  Soon we will see where this came from, but for now, observe that

 = 2xy3 + 2 = M

x

and

 = 3x2 y2 + 8e4y = N

y

With this choice of 
x y, equation (1.9) becomes



 dx + dy = 0

x

y or d
x y = 0

(1.9)

28

CHAPTER 1

First-Order Differential Equations

The general solution of this equation is 
x y = C or, in this example, x2 y3 + 2x + 2e4y = C This implicitly defines the general solution of the differential equation (1.8). To verify this, differentiate the last equation implicitly with respect to x: 2xy3 + 3x2 y2 y
+ 2 + 8e4y y
= 0 or 2xy3 + 2 +
3x2 y2 + 8e4y y
= 0 This is equivalent to the original differential equation y
= −

2xy3 + 2  3x2 y2 + 8e4y

With this as background, we will make the following definitions.

DEFINITION 1.3

Potential Function

A function  is a potential function for the differential equation M
x y + N
x yy
= 0 on a region R of the plane if, for each
x y in R,

 = M
x y and

x

DEFINITION 1.4

 = N
x y

y

Exact Differential Equation

When a potential function exists on a region R for the differential equation M + Ny
= 0, then this equation is said to be exact on R.

The differential equation of Example 1.17 is exact (over the entire plane), because we exhibited a potential function for it, defined for all
x y. Once a potential function is found, we can write an equation implicitly defining the general solution. Sometimes we can explicitly solve for the general solution, and sometimes we cannot. Now go back to Example 1.17. We want to explore how the potential function that materialized there was found. Recall we required that

 = 2xy3 + 2 = M

x

and

 = 3x2 y2 + 8e4y = N

y

1.4 Exact Differential Equations

29

Pick either of these equations to begin and integrate it. Say we begin with the first. Then integrate with respect to x: 
x y =



  dx = 2xy3 + 2 dx

x

= x2 y3 + 2x + g
y In this integration with respect to x we held y fixed, hence we must allow that y appears in the “constant” of integration. If we calculate / x, we get 2xy2 + 2 for any function g
y. Now we know  to within this function g. Use the fact that we know / y to write

 = 3x2 y2 + 8e4y

y =

2 3
x y + 2x + g
y = 3x2 y2 + g

y

y

This equation holds if g

y = 8e4y , hence we may choose g
y = 2e4y . This gives the potential function 
x y = x2 y3 + 2x + 2e4y  If we had chosen to integrate / y first, we would have gotten 
x y =




 3x2 y2 + 8e4y dy

= x2 y3 + 2e4y + h
x Here h can be any function of one variable, because no matter how h
x is chosen, 

 2 3 x y + 2e4y + h
x = 3x2 y2 + 8e4y 

y as required. Now we have two expressions for / x:

 = 2xy3 + 2

x =



 2 3 x y + 2e4y + h
x = 2xy3 + h

x

x

This equation forces us to choose h so that h

x = 2, and we may therefore set h
x = 2x. This gives 
x y = x2 y3 + 2e4y + 2x as we got before. Not every first-order differential equation is exact. For example, consider y + y
= 0

30

CHAPTER 1

First-Order Differential Equations

If there were a potential function , then we would have

 = y

x

 = 1

y

Integrate / x = y with respect to x to get 
x y = xy + g
y. Substitute this into / y = 1 to get


xy + g
y = x + g

y = 1

y But this can hold only if g

y = 1 − x, an impossibility if g is to be independent of x. Therefore, y + y
= 0 has no potential function. This differential equation is not exact (even though it is easily solved either as a separable or as a linear equation). This example suggests the need for a convenient test for exactness. This is provided by the following theorem, in which a “rectangle in the plane” refers to the set of points on or inside any rectangle having sides parallel to the axes.

THEOREM 1.1

Test for Exactness

Suppose M
x y, N
x y, M/ y, and N/ x are continuous for all
x y within a rectangle R in the plane. Then, M
x y + N
x yy
= 0 is exact on R if and only if, for each
x y in R,

M

N = 

y

x Proof

If M + Ny
= 0 is exact, then there is a potential function  and

 = M
x y and

x

 = N
x y

y

Then, for
x y in R,

M

=

y

y





x



2 

2  = = =

y x x y x





y

 =

N 

x

Conversely, suppose M/ y and N/ x are continuous on R. Choose any
x0  y0  in R and define, for
x y in R, 
x y =




x

x0

M
 y0  d +




y

N
x  d  y0

Immediately we have, from the fundamental theorem of calculus,

 = N
x y

y

(1.10)

1.4 Exact Differential Equations

31

since the first integral in equation (1.10) is independent of y. Next, compute




x


y = M
 y0  d + N
x  d

x

x x0

x y0
y N
x  d = M
x y0  + y0 x
y M = M
x y0  +
x  d y0 y = M
x y0  + M
x y − M
x y0  = M
x y and the proof is complete. For example, consider again y + y
= 0. Here M
x y = y and N
x y = 1, so

N =0

x

and

M =1

y

throughout the entire plane. Thus, y + y
= 0 cannot be exact on any rectangle in the plane. We saw this previously by showing that this differential equation can have no potential function.

EXAMPLE 1.18

Consider x2 + 3xy +
4xy + 2xy
= 0 Here M
x y = x2 + 3xy and N
x y = 4xy + 2x. Now

N = 4y + 2

x

and

M = 3x

y

and 3x = 4y + 2 is satisfied by all
x y on a straight line. However, N/ x = M/ y cannot hold for all
x y in an entire rectangle in the plane. Hence this differential equation is not exact on any rectangle.

EXAMPLE 1.19

Consider ex sin
y − 2x +
ex cos
y + 1 y
= 0 With M
x y = ex sin
y − 2x and N
x y = ex cos
y + 1, we have

M

N = ex cos
y =

x

y for all
x y. Therefore this differential equation is exact. To find a potential function, set

 = ex sin
y − 2x

x

and

 = ex cos
y + 1

y

32

CHAPTER 1

First-Order Differential Equations

Choose one of these equations and integrate it. Integrate the second equation with respect to y:

x y =
ex cos
y + 1 dy = ex sin
y + y + h
x Then we must have

 = ex sin
y − 2x

x =

x
e sin
y + y + h
x = ex sin
y + h

x

x

Then h

x = −2x and we may choose h
x = −x2 . A potential function is 
x y = ex sin
y + y − x2  The general solution of the differential equation is defined implicitly by ex sin
y + y − x2 = C Note of Caution: If  is a potential function for M + Ny
= 0,  itself is not the solution. The general solution is defined implicitly by the equation 
x y = C.

SECTION 1.4

PROBLEMS

In each of Problems 1 through 8, determine where (if anywhere) in the plane the differential equation is exact. If it is exact, find a potential function and the general solution, perhaps implicitly defined. If the equation is not exact, do not attempt a solution at this time. 1. 2y2 + yexy +
4xy + xexy + 2yy
= 0 2. 4xy + 2x +
2x2 + 3y2 y
= 0 3. 4xy + 2x2 y +
2x2 + 3y2 y
= 0 4. 2 cos
x + y − 2x sin
x + y − 2x sin
x + yy
= 0 1 5. + y +
3y2 + xy
= 0 x 6. ex sin
y2  + xex sin
y2  +
2xyex sin
y2  + ey y
= 0 7. sinh
x sinh
y + cosh
x cosh
yy
= 0 8. 4y4 + 3 cos
x +
16y3 x − 3 cos
yy
= 0 In each of Problems 9 through 14, determine if the differential equation is exact in some rectangle containing in its interior the point where the initial condition is given. If so, solve the initial value problem. This solution may be implicitly defined. If the differential equation is not exact, do not attempt a solution.

9. 3y4 − 1 + 12xy3 y
= 0 y
1 = 2 10. 2y − y2 sec2
xy2  +
2x − 2xy sec2
xy2 y
= 0 y
1 = 2 11. x cos
2y − x − sin
2y − x − 2x cos
2y − xy
= 0 y
/12 = /8 y 12. 1 + ey/x − ey/x + ey/x y
= 0 y
1 = −5 x 13. y sinh
y − x − cosh
y − x + y sinh
y − xy
= 0 y
4 = 4 14. ey +
xey − 1y
= 0 y
5 = 0 In Problems 15 and 16, choose a constant  so that the differential equation is exact, then produce a potential function and obtain the general solution. 15. 2xy3 − 3y −
3x + x2 y2 − 2yy
= 0 16. 3x2 + xy − x2 y−1 y
= 0 17. Let  be a potential function for M + Ny
= 0 in some region R of the plane. Show that for any constant c,  + c is also a potential function. How does the general solution of M + Ny
= 0 obtained by using  differ from that obtained using  + c?

1.5 Integrating Factors

1.5

33

Integrating Factors “Most” differential equations are not exact on any rectangle. But sometimes we can multiply the differential equation by a nonzero function 
x y to obtain an exact equation. Here is an example that suggests why this might be useful.

EXAMPLE 1.20

The equation y2 − 6xy +
3xy − 6x2 y
= 0

(1.11)

is not exact on any rectangle. Multiply it by 
x y = y to get y3 − 6xy2 +
3xy2 − 6x2 yy
= 0

(1.12)

Wherever y = 0, equations (1.11) and (1.12) have the same solution. The reason for this is that equation (1.12) is just

 y y2 − 6xy +
3xy − 6x2 y
= 0 and if y = 0, then necessarily y2 − 6xy +
3xy − 6x2 y
= 0. Now notice that equation (1.12) is exact (over the entire plane), having potential function 
x y = xy3 − 3x2 y2  Thus the general solution of equation (1.12) is defined implicitly by xy3 − 3x2 y2 = C and, wherever y = 0, this defines the general solution of equation (1.11) as well. To review what has just occurred, we began with a nonexact differential equation. We multiplied it by a function  chosen so that the new equation was exact. We solved this exact equation, then found that this solution also worked for the original, nonexact equation. The function  therefore enabled us to solve a nonexact equation by solving an exact one. This idea is worth pursuing, and we begin by giving a name to .

DEFINITION 1.5

Let M
x y and N
x y be defined on a region R of the plane. Then 
x y is an integrating factor for M + Ny
= 0 if 
x y = 0 for all
x y in R, and M + Ny
= 0 is exact on R.

How do we find an integrating factor for M + Ny
= 0? For  to be an integrating factor, M + Ny
= 0 must be exact (in some region of the plane), hence


N =
M

x

y

(1.13)

34

CHAPTER 1

First-Order Differential Equations

in this region. This is a starting point. Depending on M and N , we may be able to determine  from this equation. Sometimes equation (1.13) becomes simple enough to solve if we try  as a function of just x or just y.

EXAMPLE 1.21

The differential equation x − xy − y
= 0 is not exact. Here M = x − xy and N = −1 and equation (1.13) is


− =

x − xy

x

y Write this as −



 =
x − xy − x

x

y

Now observe that this equation is simplified if we try to find  as just a function of x, because in this event / y = 0 and we are left with just

 = x

x This is separable. Write 1 d = x dx  and integrate to obtain ln = 21 x2  Here we let the constant of integration be zero because we need only one integrating factor. From the last equation, choose 2


x = ex /2  a nonzero function. Multiply the original differential equation by ex
x − xye

x2 /2

−e

x2 /2

2 /2

to obtain

y
= 0

This equation is exact over the entire plane, and we find the potential function 
x y = 2
1 − yex /2 . The general solution of this exact equation is implicitly defined by
1 − yex

2 /2

= C

In this case, we can explicitly solve for y to get y
x = 1 − Ce−x /2  2

and this is also the general solution of the original equation x − xy − y
= 0. If we cannot find an integrating factor that is a function of just x or just y, then we must try something else. There is no template to follow, and often we must start with equation (1.13) and be observant.

1.5 Integrating Factors

35

EXAMPLE 1.22

Consider 2y2 − 9xy +
3xy − 6x2 y
= 0. This is not exact. With M = 2y2 − 9xy and N = 3xy − 6x2 , begin looking for an integrating factor by writing equation (1.13): 




3xy − 6x2  =


2y2 − 9xy

x

y This is
3xy − 6x2 



 + 
3y − 12x =
2y2 − 9xy + 
4y − 9x

x

y

(1.14)

If we attempt  = 
x, then / y = 0 and we obtain
3xy − 6x2 

 + 
3y − 12x = 
4y − 9x

x

which cannot be solved for  as just a function of x. Similarly, if we try  = 
y, so / x = 0, we obtain an equation we cannot solve. We must try something else. Notice that equation (1.14) involves only integer powers of x and y. This suggests that we try 
x y = xa yb . Substitute this into equation (1.14) and attempt to choose a and b. The substitution gives us 3axa yb+1 − 6axa+1 yb + 3xa yb+1 − 12xa+1 yb = 2bxa yb+1 − 9bxa+1 yb + 4xa yb+1 − 9xa+1 yb  Assume that x = 0 and y = 0. Then we can divide by xa yb to get 3ay − 6ax + 3y − 12x = 2by − 9bx + 4y − 9x Rearrange terms to write
1 + 2b − 3ay =
−3 + 9b − 6ax Since x and y are independent, this equation can hold for all x and y only if 1 + 2b − 3a = 0

and

− 3 + 9b − 6a = 0

Solve these equations to obtain a = b = 1. An integrating factor is 
x y = xy. Multiply the differential equation by xy to get 2xy3 − 9x2 y2 +
3x2 y2 − 6x3 yy
= 0 This is exact with potential function 
x y = x2 y3 − 3x3 y2 . For x = 0 and y = 0, the solution of the original differential equation is given implicitly by x2 y3 − 3x3 y2 = C The manipulations used to find an integrating factor may fail to find some solutions, as we saw with singular solutions of separable equations. Here are two examples in which this occurs.

EXAMPLE 1.23

Consider 2xy − y
= 0 y−1

(1.15)

36

CHAPTER 1

First-Order Differential Equations

We can solve this as a separable equation, but here we want to make a point about integrating factors. Equation (1.15) is not exact, but 
x y =
y − 1/y is an integrating factor for y = 0, a condition not required by the differential equation itself. Multiplying the differential equation by 
x y yields the exact equation 2x −

y−1
y = 0 y

with potential function 
x y = x2 − y + lny and general solution defined by x2 − y + lny = C

for y = 0

This is also the general solution of equation (1.15), but the method used has required that y = 0. However, we see immediately that y = 0 is also a solution of equation (1.15). This singular solution is not contained in the expression for the general solution for any choice of C.

EXAMPLE 1.24

The equation y − 3 − xy
= 0

(1.16)

is not exact, but 
x y = 1/x
y − 3 is an integrating factor for x = 0 and y = 3, conditions not required by the differential equation itself. Multiplying equation (1.16) by 
x y yields the exact equation 1 1
− y = 0 x y−3 with general solution defined by lnx + C = lny − 3 This is also the general solution of equation (1.16) in any region of the plane not containing the lines x = 0 or y = 3. This general solution can be solved for y explicitly in terms of x. First, any real number is the natural logarithm of some positive number, so write the arbitrary constant as C = ln
k, in which k can be any positive number. The equation for the general solution becomes lnx + ln
k = lny − 3 or lnkx = lny − 3 But then y − 3 = ±kx. Replacing ±k with K, which can now be any nonzero real number, we obtain y = 3 + Kx as the general solution of equation (1.16). Now observe that y = 3 is a solution of equation (1.16). This solution was “lost”, or at least not found, in using the integrating factor as a method of solution. However, y = 3 is not a singular solution because we can include it in the expression y = 3 + Kx by allowing K = 0. Thus the general solution of equation (1.16) is y = 3 + Kx, with K any real number.

1.5 Integrating Factors

1.5.1

37

Separable Equations and Integrating Factors

We will point out a connection between separable equations and integrating factors. The separable equation y
= A
xB
y is in general not exact. To see this, write it as A
xB
y − y
= 0 so in the present context we have M
x y = A
xB
y and N
x y = −1. Now


−1 = 0

x

and

A
xB
y = A
xB

y

y

and in general A
xB

y = 0. However, 
y = 1/B
y is an integrating factor for the separable equation. If we multiply the differential equation by 1/B
y, we get A
x − an exact equation because

x

−

1
y = 0 B
y

1

=

A
x = 0 B
y

y

The act of separating the variables is the same as multiplying by the integrating factor 1/B
y.

1.5.2

Linear Equations and Integrating Factors

Consider the linear equation y
+ p
xy = q
x. We can write this as p
xy − q
x + y
= 0, so in the present context, M
x y = p
xy − q
x and N
x y = 1. Now

1 = 0

x

and

p
xy − q
x = p
x

y

so the linear equation is not exact unless p
x is identically zero. However, 
x y = e is an integrating factor. Upon multiplying the linear equation by , we get

p
xy − q
xe



p
x dx

+e





p
x dx

p
x dx


y = 0

and this is exact because   

 p
x dx

 e

p
xy − q
xe p
x dx  = p
xe p
x dx =

x

y

SECTION 1.5

PROBLEMS

1. Determine a test involving M and N to tell when M + Ny
= 0 has an integrating factor that is a function of y only. 2. Determine a test to determine when M + Ny
= 0 has an integrating factor of the form 
x y = xa yb for some constants a and b. 3. Consider y − xy
= 0. (a) Show that this equation is not exact on any rectangle.

(b) Find an integrating factor 
x that is a function of x alone. (c) Find an integrating factor 
y that is a function of y alone. (d) Show that there is also an integrating factor
x y = xa yb for some constants a and b. Find all such integrating factors.

CHAPTER 1

38

First-Order Differential Equations

In each of Problems 4 through 12, (a) show that the differential equation is not exact, (b) find an integrating factor, (c) find the general solution (perhaps implicitly defined), and (d) determine any singular solutions the differential equation might have.


4. xy − 3y = 2x

13. 1 + xy
= 0 y
e4  = 0 14. 3y + 4xy
= 0 y
1 = 6 15. 2
y3 − 2 + 3xy2 y
= 0 y
3 = 1

3

16. y
1 + x + 2xy
= 0 y
4 = 6

5. 1 +
3x − e−2y y
= 0

2

6. 6x2 y + 12xy + y2 +
6x2 + 2yy
= 0

17. 2xy + 3y
= 0 y
0 = 4 (Hint: try  = ya ebx )

7. 4xy + 6y2 +
2x2 + 6xyy
= 0

18. 2y
1 + x2  + xy
= 0 y
2 = 3 (Hint: try  = xa ebx )

8. y2 + y − xy
= 0

19. sin
x − y + cos
x − y − cos
x − yy
= 0

2

y
0 = 7/6

9. 2xy2 + 2xy +
x2 y + x2 y
= 0 10. 2y2 − 9xy +
3xy − 6x2 y
= 0 (Hint: try 
x y = xa y b )


11. y + y = y (Hint: try 
x y = e y ) 4

ax b

12. x2 y
+ xy = −y−3/2 (Hint: try 
x y = xa yb ) In each of Problems 13 through 20, find an integrating factor, use it to find the general solution of the differential equation, and then obtain the solution of the initial value problem.

1.6

20. 3x2 y + y3 + 2xy2 y
= 0 y
2 = 1

21. Show that any nonzero constant multiple of an integrating factor for M + Ny
= 0 is also an integrating factor. 22. Let 
x y be an integrating factor for M + Ny
= 0 and suppose that the general solution is defined by 
x y = C. Show that 
x yG

x y is also an integrating factor, for any differentiable function G of one variable.

Homogeneous, Bernoulli, and Riccati Equations In this section we will consider three additional kinds of first-order differential equations for which techniques for finding solutions are available.

1.6.1

Homogeneous Differential Equations

DEFINITION 1.6

Homogeneous Equation

A first-order differential equation is homogeneous if it has the form y  y
= f x

1.6 Homogeneous, Bernoulli, and Riccati Equations

39

In a homogeneous equation, y
is isolated on one side, and the other side is some expression in which y and x must always appear in the combination y/x. For example, y x y
= sin y x is homogeneous, while y
= x2 y is not. Sometimes algebraic manipulation will put a first order equation into the form of the homogeneous equation. For example, y y
= (1.17) x+y is not homogeneous. However, if x = 0, we can write this as y
=

y/x  1 + y/x

(1.18)

a homogeneous equation. Any technique we develop for homogeneous equations can therefore be used on equation (1.18). However, this solution assumes that x = 0, which is not required in equation (1.17). Thus, as we have seen before, when we perform manipulations on a differential equation, we must be careful that solutions have not been overlooked. A solution of equation (1.18) will also satisfy (1.17), but equation (1.17) may have other solutions as well. Now to the point. A homogeneous equation is always transformed into a separable one by the transformation y = ux. To see this, compute y
= u
x + x
u = u
x + u and write u = y/x. Then y
= f
y/x becomes u
x + u = f
u We can write this as 1 du 1 =  f
u − u dx x or, in differential form, 1 1 du = dx f
u − u x and the variables (now x and u) have been separated. Upon integrating this equation, we obtain the general solution of the transformed equation. Substituting u = y/x then gives the general solution of the original homogeneous equation.

EXAMPLE 1.25

Consider xy
= Write this as y
=

y2 + y x

 y 2 x

y +  x

Let y = ux. Then u
x + u = u2 + u

40

CHAPTER 1

First-Order Differential Equations

or u
x = u2  Write this as 1 1 du = dx 2 u x and integrate to obtain −

1 = lnx + C u

Then u
x =

−1  lnx + C

the general solution of the transformed equation. The general solution of the original equation is y=

−x  lnx + C

EXAMPLE 1.26 A Pursuit Problem

A pursuit problem is one of determining a trajectory so that one object intercepts another. Examples involving pursuit problems are missiles fired at airplanes and a rendezvous of a shuttle with a space station. These are complex problems that require numerical approximation techniques. We will consider a simple pursuit problem that can be solved explicitly. Suppose a person jumps into a canal of constant width w and swims toward a fixed point directly opposite the point of entry into the canal. The person’s speed is v and the water current’s speed is s. Assume that, as the swimmer makes his way across, he always orients to point toward the target. We want to determine the swimmer’s trajectory. Figure 1.13 shows a coordinate system drawn so that the swimmer’s destination is the origin and the point of entry into the water is
w 0. At time t the swimmer is at the point
x
t y
t. The horizontal and vertical components of his velocity are, respectively, x

t = −v cos
 and

y

t = s − v sin


with  the angle between the positive x axis and
x
t y
t at time t. From these equations, dy y

t s − v sin
 s =
= = tan
 − sec
 dx x
t −v cos
 v

y (x, y) v

α

v sin(α) v cos(α)

(0, 0) FIGURE 1.13

(w, 0) The swimmer’s path.

x

1.6 Homogeneous, Bernoulli, and Riccati Equations

41

From Figure 1.13, tan
 =

y x

and

sec
 =

1  2 x + y2  x

Therefore y s1  2 dy = − x + y2  dx x v x Write this as the homogeneous equation y s dy = − dx x v

 1+

 y 2 x

and put y = uv to obtain s1 1 dx du = − √ 2 v x 1+u Integrate to get

   s   lnu + 1 + u2  = − lnx + C v

Take the exponential of both sides of this equation:      u + 1 + u2  = eC e−
s lnx/v  We can write this as

 u + 1 + u2 = Kx−s/v 

This equation can be solved for u. First write  1 + u2 = Kx−s/v − u and square both sides to get 1 + u2 = K 2 e−2s/v − 2Kue−s/v + u2  Now u2 cancels and we can solve for u: 1 1 s/v 1 x  u
x = Kx−s/v − 2 2K Finally, put u = y/x to get 1 1 1+s/v 1 x  y
x = Kx1−s/v − 2 2K To determine K, notice that y
w = 0, since we put the origin at the point of destination. Thus, 1 1 1 1+s/v Kw1−s/v − w =0 2 2K and we obtain K = ws/v  Therefore,



w  x 1−s/v  x 1+s/v − y
x =  2 w w

42

CHAPTER 1

First-Order Differential Equations y

s/v  34

0.30 0.25 1

0.20

s/v  2

0.15

s/v  3

0.10

s/v  5

1

1

0.05 0.0

0.2

0.4

0.6

0.8

1.0

x

Graphs of

w  x 1−s/v  x 1+s/v y= − 2 w w for s/v equal to 15 , 13 , 21 and 43 , and w chosen as 1. FIGURE 1.14

As might be expected, the path the swimmer takes depends on the width of the canal, the speed of the swimmer, and the speed of the current. Figure 1.14 shows trajectories corresponding to s/v equal to 15 , 13 , 21 and 43 , with w = 1.

1.6.2

The Bernoulli Equation

DEFINITION 1.7

A Bernoulli equation is a first order equation, y
+ P
xy = R
xy  in which  is a real number.

A Bernoulli equation is separable if  = 0 and linear if  = 1. About 1696, Leibniz showed that a Bernoulli equation with  = 1 transforms to a linear equation under the change of variables: v = y1−  This is routine to verify. Here is an example.

EXAMPLE 1.27

Consider the equation 1 y
+ y = 3x2 y3  x which is Bernoulli with P
x = 1/x, R
x = 3x2 , and  = 3. Make the change of variables v = y−2 

1.6 Homogeneous, Bernoulli, and Riccati Equations

43

Then y = v−1/2 and 1 y

x = − v−3/2 v

x 2 so the differential equation becomes 1 1 − v−3/2 v

x + v−1/2 = 3x2 v−3/2  2 x or, upon multiplying by −2v3/2 , 2 v
− v = −6x2  x a linear equation. An integrating factor is e− factor to get




2/x dx

= x−2 . Multiply the last equation by this

x−2 v
− 2x−3 v = −6 which is
x−2 v
= −6 Integrate to get x−2 v = −6x + C so v = −6x3 + Cx2  The general solution of the Bernoulli equation is y
x = 

1.6.3

1 v
x

=√

1 Cx2 − 6x3



The Riccati Equation

DEFINITION 1.8

A differential equation of the form y
= P
xy2 + Q
xy + R
x is called a Riccati equation.

A Riccati equation is linear exactly when P
x is identically zero. If we can somehow obtain one solution S
x of a Riccati equation, then the change of variables y = S
x +

1 z

44

CHAPTER 1

First-Order Differential Equations

transforms the Riccati equation to a linear equation. The strategy is to find the general solution of this linear equation and from it produce the general solution of the original Riccati equation.

EXAMPLE 1.28

Consider the Riccati equation 2 1 1 y
= y2 + y −  x x x By inspection, y = S
x = 1 is one solution. Define a new variable z by putting 1 y = 1+  z Then y
= −

1
z z2

Substitute these into the Riccati equation to get     1 2 1 1 2 1 1 1+ 1+ −  + − 2 z
= z x z x z x or 3 1 z
+ z = −  x x This is linear. An integrating factor is e




3/x dx

= x3 . Multiply by x3 to get

x3 z
+ 3x2 z =
x3 z
= −x2  Integrate to get 1 x3 z = − x3 + C 3 so 1 C z
x = − + 3  3 x The general solution of the Riccati equation is y
x = 1 +

1 1 = 1+  z
x −1/3 + C/x3

This solution can also be written y
x = in which K = 3C is an arbitrary constant.

K + 2x3  K − x3

1.6 Homogeneous, Bernoulli, and Riccati Equations

PROBLEMS

SECTION 1.6

In each of Problems 1 through 14, find the general solution. These problems include all types considered in this section. 1 1 1. y = 2 y2 − y + 1 x x


2 1 2. y + y = 3 y−4/3 x x


3. y
+ xy = xy2 4. y
=

x y + y x

5. y
=

y x+y

6. y
=

4 1 2 1 y − y− 2x x x

In each of Problems 16 through 19, use the idea of Problem 15 to find the general solution. 16. y
=

y−3 x+y−1

17. y
=

3x − y − 9 x+y+1

18. y
=

x + 2y + 7 −2x + y − 9

19. y
=

2x − 5y − 9 −4x + y + 9

20. Continuing from Problem 15, consider the case that ae − bd = 0. Now let u =
ax + by/a, assuming that a = 0. Show that this transforms the differential equation of Problem 15 into the separable equation   du b au + c = 1+ F  dx a du + r

7.
x − 2yy
= 2x − y 8. xy
= x cos
y/x + y 1 1 9. y
+ y = 4 y−3/4 x x

In each of Problems 21 through 24, use the method of Problem 20 to find the general solution.

10. x2 y
= x2 + y2 1 2 11. y
= − y2 + y x x 12. x3 y
= x2 y − y3 13. y
= −e−x y2 + y + ex 3 2 14. y
+ y = y2 x x 15. Consider the differential equation  y
= F

45

 ax + by + c  dx + ey + r

in which a, b, c, d, e, and r are constants and F is a differentiable function of one variable. (a) Show that this equation is homogeneous if and only if c = r = 0. (b) If c and/or r is not zero, this equation is called nearly homogeneous. Assuming that ae − bd = 0, show that it is possible to choose constants h and k so that the transformation X = x + h Y = y + k converts this nearly homogeneous equation into a homogeneous one. Hint: Put x = X − h y = Y − k into the differential equation and obtain a differential equation in X and Y . Use the conclusion of (a) to choose h and k so that this equation is homogeneous.

21. y
=

x−y+2 x−y+3

22. y
=

3x + y − 1 6x + 2y − 3

23. y
=

x − 2y 3x − 6y + 4

24. y
=

x−y+6 3x − 3y + 4

25. (The Pursuing Dog) A man stands at the junction of two perpendicular roads and his dog is watching him from one of the roads at a distance A feet away. At a given instant the man starts to walk with constant speed v along the other road, and at the same time the dog begins to run toward the man with speed 2v. Determine the path the dog will take, assuming that it always moves so that it is facing the man. Also determine when the dog will eventually catch the man. (This is American Mathematical Monthly problem 3942, 1941). 26. (Pursuing Bugs) One bug is located at each corner of a square table of side length a. At a given time they begin moving at constant speed v, each pursuing its neighbor to the right. (a) Determine the curve of pursuit of each bug. Hint: Use polar coordinates with the origin at the

CHAPTER 1

46

First-Order Differential Equations

center of the table and the polar axis containing one of the corners. When a bug is at
f
 , its target is at
f
  + /2. Use the chain rule to write

speed . The bug moves toward the center of the disk at constant speed v. (a) Derive a differential equation for the path of the bug, using polar coordinates.

dy dy/d =  dx dx/d

(b) How many revolutions will the disk make before the bug reaches the center? (The solution will be in terms of the angular speed and radius of the disk.)

where y
 = f
 sin
 and x
 = f
 cos
. (b) Determine the distance traveled by each bug. (c) Does any bug actually catch its quarry?

(c) Referring to (b), what is the total distance the bug will travel, taking into account the motion of the disk?

27. (The Spinning Bug) A bug steps onto the edge of a disk of radius a that is spinning at a constant angular

1.7

Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories 1.7.1

Mechanics

Before applying first-order differential equations to problems in mechanics, we will review some background. Newton’s second law of motion states that the rate of change of momentum (mass times velocity) of a body is proportional to the resultant force acting on the body. This is a vector equation, but we will for now consider only motion along a straight line. In this case Newton’s law is d F =k
mv dt We will take k = 1, consistent with certain units of measurement, such as the English, MKS, or gcs systems. The mass of a moving object need not be constant. For example, an airplane consumes fuel as it moves. If m is constant, then Newton’s law is dv = ma dt in which a is the acceleration of the object along the line of motion. If m is not constant, then F =m

dv dm +v  dt dt Newton’s law of gravitational attraction states that if two objects have masses m1 and m2 , and they (or their center of masses) are at distance r from each other, then each attracts the other with a gravitational force of magnitude mm F = G 12 2  r This force is directed along the line between the centers of mass. G is the universal gravitational constant. If one of the objects is the earth, then F =m

F =G

mM 
R + x2

where M is the mass of the earth, R is its radius (about 3,960 miles), m is the mass of the second object, and x is its distance from the surface of the earth. This assumes that the earth is

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories

47

spherical and that its center of mass is at the center of this sphere, a good enough approximation for some purposes. If x is small compared to R, then R + x is approximately R and the force on the object is approximately GM m R2 which is often written as mg. Here g = GM/R2 is approximately 32 feet per second per second or 98 meters per second per second. We are now ready to analyze some problems in mechanics. Terminal Velocity Consider an object that is falling under the influence of gravity, in a medium such as water, air or oil. This medium retards the downward motion of the object. Think, for example, of a brick dropped in a swimming pool or a ball bearing dropped in a tank of oil. We, want to analyze the object’s motion. Let v
t be the velocity at time t. The force of gravity pulls the object down and has magnitude mg. The medium retards the motion. The magnitude of this retarding force is not obvious, but experiment has shown that its magnitude is proportional to the square of the velocity. If we choose downward as the positive direction and upward as negative, then Newton’s law tells us that, for some constant , dv  dt If we assume that the object begins its motion from rest (dropped, not thrown) and if we start the clock at this instant, then v
0 = 0. We now have an initial value problem for the velocity:  v
= g − v2  v
0 = 0 m This differential equation is separable. In differential form, F = mg − v2 = m

1 dv = dt g −
/mv2 Integrate to get



m tanh−1 g

Solve for the velocity, obtaining v
t =





mg tanh 

  v = t + C mg



 g
t + C  m

Now use the initial condition to solve for the integration constant:     mg g v
0 = tanh C = 0  m Since tanh
 = 0 only if  = 0, this requires that C = 0 and the solution for the velocity is    mg g tanh t  v
t =  m Even in this √ generality, we can draw an important conclusion about the motion. As t increases, tanh
g/mt approaches 1. This means that  mg  lim v
t = t→ 

CHAPTER 1

48

First-Order Differential Equations

This means that an object falling under the influence of gravity, through a retarding medium (with force proportional to the square of the velocity), will not√ increase in velocity indefinitely. Instead, the object’s velocity approaches the limiting value mg/. If the medium is deep enough, the object will settle into a descent of approximately constant velocity. This number √ mg/ is called the terminal velocity of the object. Skydivers experience this phenomenon.

7 ft 9 ft

Motion of a Chain on a Pulley A 16 foot long chain weighing  pounds per foot hangs over a small pulley, which is 20 feet above the floor. Initially, the chain is held at rest with 7 feet on one side and 9 on the other, as in Figure 1.15. How long after the chain is released, and with what velocity, will it leave the pulley? When 8 feet of chain hang on each side of the pulley, the chain is in equilibrium. Call this position x = 0 and let x
t be the distance the chain has fallen below this point at time t. The net force acting on the chain is 2x and the mass of the chain is 16/32, or /2 slugs. The ends of the chain have the same speed as its center of mass, so the acceleration of the chain at its center of mass is the same as it is at its ends. The equation of motion is  dv = 2x 2 dt

FIGURE 1.15

Chain on a pulley.

from which  cancels to yield dv = 4x dt A chain rule differentiation enables us to write this equation in terms of v as a function of x. Write dv dv dx dv = =v  dt dx dt dx Then v

dv = 4x dx

This is a separable equation, which we solve to get v2 = 4x2 + K Now, x = 1 when v = 0, so K = −4 and v2 = 4x2 − 4 The chain leaves the pulley when x = 8. Whenever this occurs, v2 = 4
63 = 252, so v = feet per second (about 1587 feet per second). To calculate the time tf required for the chain to leave the pulley, compute tf = =




tf 0




8 1




√ 6 7

dt dv dv 0
81 dt dx = dv dx 1 v dt =

√

√ 252 = 6 7

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories √ Since v
x = 2 x2 − 1,

49

 8 √ 1  1 1
8  2 lnx + x − 1 dx = tf = √ 2 1 2 x2 − 1 1 √ 1 = ln
8 + 63 2

about 138 seconds. In this example the mass was constant, so dm/dt = 0 in Newton’s law of motion. Next is an example in which the mass varies with time. Chain Piling on the Floor Suppose a 40 foot long chain weighing  pounds per foot is supported in a pile several feet above the floor, and begins to unwind when released from rest with 10 feet already played out. Determine the velocity with which the chain leaves the support. The amount of chain that is actually in motion changes with time. Let x
t denote the length of that part of the chain that has left the support by time t and is currently in motion. The equation of motion is dm dv +v = F (1.19) dt dt where F is the total external force acting on the chain. Now F = x = mg, so m = x/g = x/32. Then m

dm  dx  = = v dt 32 dt 32 Further, dv dv =v  dt dx as in the preceding example. Put this information into equation (1.19) to get  x dv v + v2 = x 32 dx 32 If we multiply this equation by 32/xv, we get 32 dv 1 + v=  dx x v

(1.20)

which we recognize as a Bernoulli equation with  = −1. Make the transformation w = v1− = v2 . Then v = w1/2 and dv 1 −1/2 dw  = w dx 2 dx Substitute these into equation (1.20) to get 1 −1/2 dw 1 1/2 w + w = 32w−1/2  2 dx x Upon multiplying this equation by 2w1/2 , we get w
+

2 w = 64 x

a linear equation for w
x. Solve this to get w
x = v
x2 =

64 C x+ 2 3 x

50

CHAPTER 1

First-Order Differential Equations

Since v = 0 when x = 10, 0 =
64/3
10 + C/100, so C = −64 000/3. Therefore,

64 1000 2 v
x = x− 2  3 x The chain leaves the support when x = 40. At this time,

64 1000 2 v = 40 − = 4
210 3 1600 √ so, the velocity is v = 2 210, or about 29 feet per second. In these models involving chains, air resistance was neglected as having no significant impact on the outcome. This was quite different from the analysis of terminal velocity, in which air resistance is a key factor. Without it, skydivers dive only once! Motion of a Block Sliding on an Inclined Plane A block weighing 96 pounds is released from rest at the top of an inclined plane of slope length 50 feet, √ and making an angle /6 radians with the horizontal. Assume a coefficient of friction of  = 3/4. Assume also that air resistance acts to retard the block’s descent down the ramp, with a force of magnitude equal to one half the block’s velocity. We want to determine the velocity v
t of the block at any time t. Figure 1.16 shows the forces acting on the block. Gravity acts downward with magnitude mg sin
, which is 96 sin
/6, or 48 pounds. Here mg = 96 is the weight of the block. The drag due to friction acts in the reverse direction and is, in pounds, √  3 −N = −mg cos
 = −
96 cos = −36 4 6 The drag force due to air resistance is −v/2, the negative sign indicating that this is a retarding force. The total external force on the block is F = 48 − 36 − 21 v = 12 − 21 v Since the block weighs 96 pounds, it has a mass of 96/32 slugs, or 3 slugs. From Newton’s second law, 3

dv 1 = 12 − v dt 2

This is a linear equation, which we write as 1 v
+ v = 4 6 N mg sin(θ) θ

mg cos(θ)

mg π 6

FIGURE 1.16

inclined plane.

Forces acting on a block on an

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories An integrating factor is e




1/6 dt

51

= et/6 . Multiply the differential equation by this factor to obtain 
 v
et/6 + 16 et/6 v = vet/6 = 4et/6

and integrate to get vet/6 = 24et/6 + C The velocity is v
t = 24 + Ce−t/6  Since the block starts from rest at time zero, v
0 = 0 = 24 + C, so C = −24 and   v
t = 24 1 − e−t/6  Let x
t be the position of the block at any time, measured from the top of the plane. Since v
t = x

t, we get
x
t = v
t dt = 24t + 144e−t/6 + K If we let the top of the block be the origin along the inclined plane, then x
0 = 0 = 144 + K, so K = −144 The position function is

  x
t = 24t + 144 e−t/6 − 1 

We can now determine the block’s position and velocity at any time. Suppose, for example, we want to know when the block reaches the bottom of the ramp. This happens when the block has gone 50 feet. If this occurs at time T , then   x
T = 50 = 24T + 144 e−T/6 − 1  This transcendental equation cannot be solved algebraically for T , but a computer approximation yields T ≈ 58 seconds. Notice that lim v
t = 24

t→

which means that the block sliding down the ramp has a terminal velocity. If the ramp is long enough, the block will eventually settle into a slide of approximately constant velocity. The mathematical model we have constructed for the sliding block can be used to analyze the motion of the block under a variety of conditions. For example, we can solve the equations leaving  arbitrary, and determine the influence of the slope angle of the ramp on position and velocity. Or we could leave  unspecified and study the influence of friction on the motion.

1.7.2

Electrical Circuits

Electrical engineers often use differential equations to model circuits. The mathematical model is used to analyze the behavior of circuits under various conditions, and aids in the design of circuits having specific characteristics. We will look at simple circuits having only resistors, inductors and capacitors. A capacitor is a storage device consisting of two plates of conducting material isolated from one another by an insulating material, or dielectric. Electrons can be transferred from one plate to another via external circuitry by applying an electromotive force to the circuit. The charge on a capacitor is essentially a count of the difference between the numbers of electrons on the two plates. This charge is proportional to the applied electromotive force, and the constant of proportionality

52

CHAPTER 1

First-Order Differential Equations

is the capacitance. Capacitance is usually a very small number, given in micro (10−6 ) or pico (10−12 ) farads. To simplify examples and problems, some of the capacitors in this book are assigned numerical values that would actually make them occupy large buildings. An inductor is made by winding a conductor such as wire around a core of magnetic material. When a current is passed through the wire, a magnetic field is created in the core and around the inductor. The voltage drop across an inductor is proportional to the change in the current flow, and this constant of proportionality is the inductance of the inductor, measured in henrys. Current is measured in amperes, with one amp equivalent to a rate of electron flow of one coulomb per second. Charge q
t and current i
t are related by i
t = q

t The voltage drop across a resistor having resistance R is iR. The drop across a capacitor having capacitance C is q/C. And the voltage drop across an inductor having inductance L is Li

t. We construct equations for a circuit by using Kirchhoff’s current and voltage laws. Kirchhoff’s current law states that the algebraic sum of the currents at any juncture of a circuit is zero. This means that the total current entering the junction must balance the current leaving (conservation of energy). Kirchhoff’s voltage law states that the algebraic sum of the potential rises and drops around any closed loop in a circuit is zero. As an example of modeling a circuit mathematically, consider the circuit of Figure 1.17. Starting at point A, move clockwise around the circuit, first crossing the battery, where there is an increase in potential of E volts. Next there is a decrease in potential of iR volts across the resistor. Finally, there is a decrease of Li

t across the inductor, after which we return to point A. By Kirchhoff’s voltage law, E − iR − Li
= 0 which is the linear equation i
+

E E i=  R L

Solve this to obtain E + Ke−Rt/L  R To determine the constant K, we need to be given the current at some time. Even without this, we can tell from this equation that as t → , the current approaches the limiting value E/R. This is the steady-state value of the current in the circuit. Another way to derive the differential equation of this circuit is to designate one of the components as a source, then set the voltage drop across that component equal to the sum of the voltage drops across the other components. To see this approach, consider the circuit of Figure 1.18. Suppose the switch is initially open so that no current flows, and that the charge i
t =

R

R L

E

C E

A FIGURE 1.17

RL Circuit.

FIGURE 1.18

RC circuit.

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories

53

on the capacitor is zero. At time zero, close the switch. We want the charge on the capacitor. Notice that we have to close the switch before there is a loop. Using the battery as a source, write iR +

1 q = E C

or Rq
+

1 q = E C

This leads to the linear equation q
+

1 E q=  RC R

with solution   q
t = EC 1 − e−t/RC satisfying q
0 = 0. This equation provides a good deal information about the circuit. Since the voltage on the capacitor at time t is q
t/C, or E
1 − e−t/RC , we can see that the voltage approaches E as t → . Since E is the battery potential, the difference between battery and capacitor voltages becomes negligible as time increases, indicating a very small voltage drop across the resistor. The current in this circuit can be computed as i
t = q

t =

E −t/RC e R

after the circuit is switched on. Thus i
t → E/R as t → . Often we encounter discontinuous currents and potential functions in dealing with circuits. These can be treated using Laplace transform techniques, which we will discuss in Chapter 3.

1.7.3

Orthogonal Trajectories

Two curves intersecting at a point P are said to be orthogonal if their tangents are perpendicular (orthogonal) at P. Two families of curves, or trajectories, are orthogonal if each curve of the first family is orthogonal to each curve of the second family, wherever an intersection occurs. Orthogonal families occur in many contexts. Parallels and meridians on a globe are orthogonal, as are equipotential and electric lines of force. A problem that occupied Newton and other early developers of the calculus was the determination of the family of orthogonal trajectories of a given family of curves. Suppose we are given a family
of curves in the plane. We want to construct a second family  of curves so that every curve in
is orthogonal to every curve in  wherever an intersection occurs. As a simple example, suppose
consists of all circles about the origin. Then  consists of all straight lines through the origin (Figure 1.19). It is clear that each straight line is orthogonal to each circle wherever the two intersect.

54

CHAPTER 1

First-Order Differential Equations y

x

FIGURE 1.19 Orthogonal families: circles and lines.

In general, suppose we are given a family
of curves. These must be described in some way, say by an equation F
x y k = 0 giving a different curve for each choice of the constant k. Think of these curves as integral curves of a differential equation y
= f
x y which we determine from the equation F
x y k = 0 by differentiation. At a point
x0  y0 , the slope of the curve C in
through this point is f
x0  y0 . Assuming that this is nonzero, any curve through
x0  y0  and orthogonal to C at this point, must have slope −1/f
x0  y0 . (Here we use the fact that two lines are orthogonal if and only if their slopes are negative reciprocals.) The family  of orthogonal trajectories of
therefore consists of the integral curves of the differential equation y
= −

1  f
x y

Solve this differential equation for the curves in .

EXAMPLE 1.29

Consider the family
of curves that are graphs of F
x y k = y − kx2 = 0 This is a family of parabolas. We want the family of orthogonal trajectories. First obtain the differential equation of
. Differentiate y − kx2 = 0 to get y
− 2kx = 0 To eliminate k, use the equation y − kx2 = 0 to write y k= 2 x Then y y
− 2 2 x = 0 x or y y
= 2 = f
x y x

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories

55

This is the differential equation of the family
. Curves in
are integral curves of this differential equation, which is of the form y
= f
x y, with f
x y = 2y/x. The family  of orthogonal trajectories therefore has differential equation y
= −

1 x =−  f
x y 2y

This equation is separable, since 2ydy = −x dx Integrate to get y2 = − 21 x2 + C This is a family of ellipses 1 2 x + y2 2

= C

Some of the parabolas and ellipses from
and  are shown in Figure 1.20. Each parabola in
is orthogonal to each ellipse in  wherever these curves intersect. y

x

FIGURE 1.20 Orthogonal families: parabolas and ellipses.

SECTION 1.7

PROBLEMS

Mechanical Systems 1. Suppose that the pulley described in this section is only 9 feet above the floor. Assuming the same initial conditions as in the discussion, determine the velocity with which the chain leaves the pulley. Hint: The mass of the part of the chain that is in motion is
16 − x/32. 2. Determine the time it takes for the chain in Problem 1 to leave the pulley.

 pounds per foot. The chain is supported by a small frictionless pulley located more than 40 feet above the floor. Initially, the chain is held at rest with 23 feet hanging on one side of the pulley with the remainder of the chain, along with the weight, on the other side. How long after the chain is released, and with what velocity, will it leave the pulley?

3. Suppose the support is only 10 feet above the floor in the discussion of the chain piling on the floor. Calculate the velocity of the moving part of the chain as it leaves the support. (Note the hint to Problem 1.)

5. (Chain on a Table) A 24-foot chain weighing  pounds per foot is stretched out on a very tall, frictionless table with 6 feet hanging off the edge. If the chain is released from rest, determine the time it takes for the end of the chain to fall off the table, and also the velocity of the chain at this instant.

4. (Chain and Weight on a Pulley) An 8-pound weight is attached to one end of a 40-foot chain that weighs

6. (Variable Mass Chain on a Low Table) Suppose the chain in Problem 5 is placed on a table that is only

56

CHAPTER 1

First-Order Differential Equations

4 feet high, so that the chain accumulates on the floor as it slides off the table. Two feet of chain are already piled up on the floor at the time that the rest of the chain is released. Determine the velocity of the moving end of the chain at the instant it leaves the table top. Hint: The mass of that part of the chain that is moving changes with time. Newton’s law applies to the center of mass of the moving system. 7. Determine the time it takes for the chain to leave the support in the discussion of the chain piling on the floor. 8. Use the conservation of energy principle (potential energy plus kinetic energy of a conservative system is a constant of the motion) to obtain the velocity of the chain in the discussion involving the chain on the pulley. 9. Use the conservation of energy principle to give an alternate derivation of the conclusion of the discussion of the chain piling on the floor. 10. (Paraboloid of Revolution) Determine the shape assumed by the surface of a liquid being spun in a circular bowl at constant angular velocity . Hint: Consider a particle of liquid located at
x y on the surface of the liquid, as in Figure 1.21. The forces acting on the particle are the horizontal force having magnitude m2 x and a vertical force of magnitude mg. Since the particle is in radial equilibrium, the resultant vector is normal to the curve.

y (x, y)

mω2 x

mg x FIGURE 1.21 Particle on the surface of a spinning liquid.

Properties of spinning liquids have found application in astronomy. A Canadian astronomer has constructed a telescope by spinning a bowl of mercury, creating a reflective surface free of the defects obtained by the usual grinding of a solid lens. He claims that the idea was probably known to Newton, but that he is the first to carry it out in practice. Roger Angel, a University of Arizona astronomer, has developed this idea into a technique for producing telescope mirrors called spin casting. As reported in Time (April 27, 1992), “   a complex ceramic mold is assembled inside the furnace and filled with glittering chunks of Pyrex-type glass. Once the furnace lid is sealed,

the temperature will slowly ratchet up over a period of several days, at times rising no more than 2 degrees Centigrade in an hour. At 750 degrees C (1382 degrees Fahrenheit), when the glass is a smooth, shiny lake, the furnace starts to whirl like a merry-go-round, an innovation that automatically spins the glass into the parabolic shape traditionally achieved by grinding.” The result is a parabolic surface requiring little or no grinding before a reflective coat is applied. Professor Angel believes that the method will allow the construction of much larger mirrors than are possible by conventional techniques. Supporting this claim is his recent production of one of the world’s largest telescope mirrors, a 65 meter (about 21 feet) to be placed in an observatory atop Mount Hopkins in Arizona. 11. A 10-pound ballast bag is dropped from a hot air balloon which is at an altitude of 342 feet and ascending at a rate of 4 feet per second. Assuming that air resistance is not a factor, determine the maximum height attained by the bag, how long it remains aloft, and the speed with which it strikes the ground. 12. A 48-pound box is given an initial push of 16 feet per second down an inclined plane that has a gradient of 7 . If there is a coefficient of friction of 13 between 24 the box and the plane, and an air resistance equal to 3 the velocity of the box, determine how far the box 2 will travel before coming to rest. 13. A skydiver and her equipment together weigh 192 pounds. Before the parachute is opened, there is an air drag equal to six times her velocity. Four seconds after stepping from the plane, the skydiver opens the parachute, producing a drag equal to three times the square of the velocity. Determine the velocity and how far the skydiver has fallen at time t. What is the terminal velocity? 14. Archimedes’ principle of buoyancy states that an object submerged in a fluid is buoyed up by a force equal to the weight of the fluid that is displaced by the object. A rectangular box, 1 by 2 by 3 feet, and weighing 384 pounds, is dropped into a 100-foot-deep freshwater lake. The box begins to sink with a drag due to the water having magnitude equal to 21 the velocity. Calculate the terminal velocity of the box. Will the box have achieved a velocity of 10 feet per second by the time it reaches bottom? Assume that the density of the water is 625 pounds per cubic foot. 15. Suppose the box in Problem 14 cracks open upon hitting the bottom of the lake, and 32 pounds of its contents fall out. Approximate the velocity with which the box surfaces. 16. The acceleration due to gravity inside the earth is proportional to the distance from the center of the

1.7 Applications to Mechanics, Electrical Circuits, and Orthogonal Trajectories earth. An object is dropped from the surface of the earth into a hole extending through the earth’s center. Calculate the speed the object achieves by the time it reaches the center. 17. A particle starts from rest at the highest point of a vertical circle and slides under only the influence of gravity along a chord to another point on the circle. Show that the time taken is independent of the choice of the terminal point. What is this common time?

57

10 

30 

15 

i1

i2

i3

5H

6V

1 10

F

4 10

F

FIGURE 1.24

Circuits 18. Determine each of the currents in the circuit of Figure 1.22.

10 

22. In a constant electromotive force RL circuit, we find that the current is given by  E i
t = 1 − e−Rt/L + i
0e−Rt/L  R Let i
0 = 0. (a) Show that the current increases with time. (b) Find a time t0 at which the current is 63% of E/R. This time is called the inductive time constant of the circuit.

i2

i1 15 

10 V

30 

FIGURE 1.22

19. In the circuit of Figure 1.23, the capacitor is initially discharged. How long after the switch is closed will the capacitor voltage be 76 volts? Determine the current in the resistor at that time. (Here k denotes 1000 ohms and F denotes 10−6 farads.)

(c) Does the inductive time constant depend on i
0? If so, in what way? 23. Recall that the charge q
t in an RC circuit satisfies the linear differential equation 1 1 q = E
t q
+ RC R (a) Solve for the charge in the case that E
t = E, constant. Evaluate the constant of integration by using the condition q
0 = q0 . (b) Determine limt→ q
t and show that this limit is independent of q0 . (c) Graph q
t. Determine when the charge has its maximum and minimum values.

250  2 µF 80 V FIGURE 1.23

20. Suppose, in Problem 19, the capacitor had a potential of 50 volts when the switch was closed. How long would it take for the capacitor voltage to reach 76 volts? 21. For the circuit of Figure 1.24, find all currents immediately after the switch is closed, assuming that all of these currents and the charges on the capacitors are zero just prior to closing the switch.

(d) Determine at what time q
t is within 1% of its steady-state value (the limiting value requested in (b)).

Orthogonal Trajectories In each of Problems 24 through 29, find the family of orthogonal trajectories of the given family of curves. If software is available, graph some curves in the given family and some curves in the family of orthogonal trajectories. 24. x + 2y = K 25. 2x2 − 3y = K 26. x2 + 2y2 = K 27. y = Kx2 + 1 28. x2 − Ky2 = 1 29. y = ekx

58

1.8

CHAPTER 1

First-Order Differential Equations

Existence and Uniqueness for Solutions of Initial Value Problems We have solved several initial value problems y
= f
x y

y
x0  = y0 

and have always found that there is just one solution. That is, the solution existed, and it was unique. Can either existence or uniqueness fail to occur? The answer is yes, as the following examples show.

EXAMPLE 1.30

Consider the initial value problem y
= 2y1/2 

y
0 = −1

The differential equation is separable and has general solution y
x =
x + C2  To satisfy the initial condition, we must choose C so that y
0 = C 2 = −1 and this is impossible if C is to be a real number. This initial value problem has no real-valued solution.

EXAMPLE 1.31

Consider the problem y
= 2y1/2 

y
2 = 0

One solution is the trivial function y = 
x = 0 But there is another solution. Define 
x =



0
x − 22

for all x

for x ≤ 2 for x ≥ 2

Graphs of both solutions are shown in Figure 1.25. Uniqueness fails in this example. Because of examples such as these, we look for conditions that ensure that an initial value problem has a unique solution. The following theorem provides a convenient set of conditions.

1.8 Existence and Uniqueness for Solutions of Initial Value Problems y

59

y y  (x)

y  (x)

x

x 2 √

FIGURE 1.25 Graphs of solutions of y
= 2 y;

y
2 = 0.

THEOREM 1.2

Existence and Uniqueness

Let f and f/ y be continuous for all
x y in a closed rectangle R centered at
x0  y0 . Then there exists a positive number h such that the initial value problem y
= f
x y

y
x0  = y0

has a unique solution defined over the interval
x0 − h x0 + h. As with the test for exactness (Theorem 1.1), by a closed rectangle we mean all points on or inside a rectangle in the plane, having sides parallel to the axes. Geometrically, existence of a solution of the initial value problem means that there is an integral curve of the differential equation passing through
x0  y0 . Uniqueness means that there is only one such curve. This is an example of a local theorem, in the following sense. The theorem guarantees existence of a unique solution that is defined on some interval of width 2h, but it says nothing about how large h is. Depending on f and x0 , h may be small, giving us existence and uniqueness “near” x0 . This is dramatically demonstrated by the initial value problem y
= y2 

y
0 = n

in which n is any positive integer. Here f
x y = y2 and f/ y = 2y, both continuous over the entire plane, hence on any closed rectangle about
0 n. The theorem tells us that there is a unique solution of this initial value problem in some interval
−h h. In this case we can solve the initial value problem explicitly, obtaining y
x = −

1  x − n1

This solution is valid for −1/n < x < 1/n, so we can take h = 1/n in this example. This means that the size of n in the initial value controls the size of the interval for the solution. The larger n is, the smaller this interval must be. This fact is certainly not apparent from the initial value problem itself! In the special case that the differential equation is linear, we can improve considerably on the existence/uniqueness theorem. THEOREM 1.3

Let p and q be continuous on an open interval I and let x0 be in I. Let y0 be any number. Then the initial value problem y
+ p
xy = q
x has a unique solution defined for all x in I.

y
x0  = y0

60

CHAPTER 1

First-Order Differential Equations

In particular, if p and q are continuous for all x, then there is a unique solution defined over the entire real line. Equation (1.6) of Section 1.3 gives the general solution of the linear equation. Using this, we can write the solution of the initial value problem:
x

x x − x p
 d p
 d x q
e 0 d + y0  y
x = e 0

Proof

x0

Because p and q are continuous on I, this solution is defined for all x in I. Therefore, in the case that the differential equation is linear, the initial value problem has a unique solution in the largest open interval containing x0 , in which both p and q are continuous.

PROBLEMS

SECTION 1.8

In each of Problems 1 through 5, show that the conditions of Theorem 1.2 are satisfied by the initial value problem. Assume familiar facts from the calculus about continuity of real functions of one and two variables. 1. y
= 2y2 + 3xey sin
xy y
2 = 4 2. y
= 4xy + cosh
x y
1 = −1 3. y
=
xy3 − sin
y y
2 = 2 4. y
= x5 − y5 + 2xey  y
3 =  5. y
= x2 ye−2x + y2  y
3 = 8 6. Consider the initial value problem y
 = 2y y
x0  = y0 . (a) Find two solutions, assuming that y0 > 0. (b) Explain why part (a) does not violate Theorem 1.2. Theorem 1.2 can be proved using Picard iterates, which we will discuss briefly. Suppose f and f/ y are continuous in a closed rectangle R having
x0  y0  in its interior and sides parallel to the axes. Consider the initial value problem y
= f
x y y
x0  = y0 . For each positive integer n, define
x yn
x = y0 + f
t yn−1
t dt x0

This is a recursive definition, giving y1
x in terms of y0 , then y2
x in terms of y1
x, and so on. The functions yn
x for n = 1 2    are called Picard iterates for the initial value problem. Under the assumptions made on f , the sequence yn
x converges for all x in some interval about x0 , and the limit of this sequence is the solution of the initial value problem on this interval. In each of Problems 7 through 10, (a) use Theorem 1.2 to show that the problem has a solution in some interval about x0 , (b) find this solution, (c) compute Picard iterates y1
x through y6
x, and from these guess yn
x in general, and (d) find the Taylor series of the solution from (b) about x0 . You should find that the iterates computed in (c) are partial sums of the series of (d). Conclude that in these examples the Picard iterates converge to the solution. 7. y
= 2 − y y
0 = 1 8. y
= 4 + y y
0 = 3 9. y
= 2x2  y
1 = 3 10. y
= cos
x y
 = 1

2

CHAPTER

REDUCTION OF ORDER CONSTANT COEFFICIENT HOMOGENEOUS LINEAR EQUATION EULER’S EQUATION NONHOMOGENEOUS EQUATION y

+ p
xy
+ q
xy = f
x APPLICATION OF SECOND ORDER DIFFERENTIAL

Second-Order Differential Equations

2.1

Preliminary Concepts A second-order differential equation is an equation that contains a second derivative, but no higher derivative. Most generally, it has the form F
x y y
 y

 = 0 although only a term involving y

need appear explicitly. For example, y

= x3  xy

− cos
y = ex and y

− 4xy
+ y = 2 are second-order differential equations. A solution of F
x y y
 y

 = 0 on an interval I (perhaps the whole real line) is a function  that satisfies the differential equation at each point of I: F
x 
x 

x 


x = 0

for x in I.

For example, 
x = 6 cos
4x − 17 sin
4x is a solution of y

+ 16y = 0 for all real x. And 
x = x3 cos
ln
x is a solution of x2 y

− 5xy
+ 10y = 0 for x > 0. These can be checked by substitution into the differential equation. The linear second-order differential equation has the form R
xy

+ P
xy
+ Q
xy = S
x 61

62

CHAPTER 2

Second-Order Differential Equations

in which R, P, Q, and S are continuous in some interval. On any interval where R
x = 0, we can divide this equation by R
x and obtain the special linear equation y

+ p
xy
+ q
xy = f
x

(2.1)

For the remainder of this chapter, we will concentrate on this equation. We want to know: 1. What can we expect in the way of existence and uniqueness of solutions of equation (2.1)? 2. How can we produce all solutions of equation (2.1), at least in some cases that occur frequently and have important applications? We begin with the underlying theory that will guide us in developing techniques for explicitly producing solutions of equation (2.1).

2.2

Theory of Solutions of y

+ p
xy
+ q
xy = f
x To get some feeling for what we are dealing with, and what we should be looking for, consider the simple linear second-order equation y

− 12x = 0 We can write this as y

= 12x and integrate to obtain y
=




y


x dx =




12x dx = 6x2 + C

Integrate again: y
x =




y

x dx =





6x2 + C dx = 2x3 + Cx + K

This solution is defined for all x, and contains two arbitrary constants. If we recall that the general solution of a first order equation contained one arbitrary constant, it seems natural that the solution of a second-order equation, involving two integrations, should contain two arbitrary constants. For any choices of C and K, we can graph the integral curves y = 2x3 + Cx + K as curves in the plane. Figure 2.1 shows some of these curves for different choices of these constants. Unlike the first-order case, there are many integral curves through each point in the plane. For example, suppose we want a solution satisfying the initial condition y
0 = 3 Then we need y
0 = K = 3 but are still free to choose C as any number. All solutions y
x = 2x3 + Cx + 3 pass through
0 3. Some of these curves are shown in Figure 2.2.

2.2 Theory of Solutions of y

+ p
xy
+ q
xy = f
x

63

y 20

10


2


1

x

1

2


10


20

FIGURE 2.1 Graphs of y = 2x3 + Cx + K for various values of C and K.

y 30 20 10
2


1 0

1

2

x


10
20 FIGURE 2.2 Graphs of y = 2x3 + Cx + 3 for various values of C.

We single out exactly one of these curves if we specify its slope at
0 3. Suppose, for example, we also specify the initial condition y

0 = −1 Since y

x = 6x2 + C, this requires that C = −1. There is exactly one solution satisfying both initial conditions (going through a given point with given slope), and it is y
x = 6x2 − x + 3 A graph of this solution is given in Figure 2.3. To sum up, at least in this example, the general solution of the differential equation involved two arbitrary constants. An initial condition y
0 = 3, specifying that the solution curve must pass through
0 3, determined one of these constants. However, that left infinitely many

64

CHAPTER 2

Second-Order Differential Equations y 40 20


2
3

0
1
20

1

2

3

x


40 FIGURE 2.3

Graph of y = 2x3 − x + 3.

solution curves passing through
0 3. The other initial condition, y

0 = −1, picked out that solution curve through
0 3 having slope −1 and gave a unique solution of this problem. This suggests that we define the initial value problem for equation (2.1) to be the differential equation, defined on some interval, together with two initial conditions, one specifying a point lying on the solution curve, and the other its slope at that point. This problem has the form y

+ p
xy
+ q
xy = f
x y
x0  = A y

x0  = B in which A and B are given real numbers. The main theorem on existence and uniqueness of solutions for this problem is the secondorder analogue of Theorem 1.3 in Chapter 1. THEOREM 2.1

Let p, q, and f be continuous on an open interval I. Let x0 be in I and let A and B be any real numbers. Then the initial value problem y

+ p
xy
+ q
xy = f
x y
x0  = A y

x0  = B has a unique solution defined for all x in I. This gives us an idea of the kind of information needed to specify a unique solution of equation (2.1). Now we need a framework in which to proceed in finding solutions. We will provide this in two steps, beginning with the case that f
x is identically zero.

2.2.1

The Homogeneous Equation y

+ p
xy
+ q
x = 0

When f
x is identically zero in equation (2.1), the resulting equation y

+ p
xy
+ q
x = 0

(2.2)

is called homogeneous. This term was used in a different context with first-order equations, and its use here is unrelated to that. Here it simply means that the right side of equation (2.1) is zero. A linear combination of solutions y1
x and y2
x of equation (2.2) is a sum of constant multiples of these functions: c1 y1
x + c2 y2
x with c1 and c2 real numbers. It is an important property of the homogeneous linear equation that linear combinations of solutions are again solutions.

2.2 Theory of Solutions of y

+ p
xy
+ q
xy = f
x

65

THEOREM 2.2

Let y1 and y2 be solutions of y

+ p
xy
+ q
xy = 0 on an interval I. Then any linear combination of these solutions is also a solution. Let c1 and c2 be real numbers. Substituting y
x = c1 y1
x + c2 y2
x into the differential equation, we obtain

Proof


c1 y1 + c2 y2 

+ p
x
c1 y1 + c2 y2 
+ q
x
c1 y1 + c2 y2  = c1 y1

+ c2 y2

+ c1 p
xy1
+ c2 p
xy2
+ c1 q
xy1 + c2 q
xy2 = c1 y1

+ p
xy1
+ q
xy1  + c2 y2

+ p
xy2
+ q
xy2  = 0 + 0 = 0 because of the assumption that y1 and y2 are both solutions. Of course, as a special case (c2 = 0), this theorem tells us also that, for the homogeneous equation, a constant multiple of a solution is a solution. Even this special case of the theorem fails for a nonhomogeneous equation. For example, y1
x = 4e2x /5 is a solution of y

+ 2y
− 3y = 4e2x  but 5y1
x = 4e2x is not. The point to taking linear combinations c1 y1 + c2 y2 is to obtain more solutions from just two solutions of equation (2.2). However, if y2 is already a constant multiple of y1 , then c1 y1 + c2 y2 = c1 y1 + c2 ky1 =
c1 + kc2 y1  just another constant multiple of y1 . In this event, y2 is superfluous, providing us nothing we did not know from just y1 . This leads us to distinguish the case in which one solution is a constant multiple of another, from the case in which the two solutions are not multiples of each other.

DEFINITION 2.1

Linear Dependence, Independence

Two functions f and g are linearly dependent on an open interval I if, for some constant c, either f
x = cg
x for all x in I, or g
x = cf
x for all x in I. If f and g are not linearly dependent on I, then they are said to be linearly independent on the interval.

EXAMPLE 2.1

y1
x = cos
x and y2
x = sin
x are solutions of y

+ y = 0, over the real line. Neither of these functions is a constant multiple of the other. Indeed, if cos
x = k sin
x for all x, then in particular √    √2 2 = = k sin =k  cos 4 2 4 2 so k must be 1. But then cos
x = sin
x for all x, a clear absurdity (for example, let x = 0). These solutions are linearly independent. Now we know from Theorem 2.2 that a cos
x + b sin
x

66

CHAPTER 2

Second-Order Differential Equations

is a solution for any numbers a and b. Because cos
x and sin
x are linearly independent, this linear combination provides an infinity of new solutions, instead of just constant multiples of one we already know. There is a simple test to tell whether two solutions of equation (2.2) are linearly independent on an interval. Define the Wronskian of solutions y1 and y2 to be W
x = y1
xy2

x − y1

xy2
x This is the 2 × 2 determinant

THEOREM 2.3

 y
x  1 W
x = 
y1
x

 y2
x  y2

x

Wronskian Test

Let y1 and y2 be solutions of y

+ p
xy
+ q
xy = 0 on an open interval I. Then, 1. Either W
x = 0 for all x in I, or W
x = 0 for all x in I. 2. y1 and y2 are linearly independent on I if and only if W
x = 0 on I. Conclusion (1) means that the Wronskian of two solutions cannot be nonzero at some points of I and zero at others. Either the Wronskian vanishes over the entire interval, or it is nonzero at every point of the interval. Conclusion (2) states that nonvanishing of the Wronskian is equivalent to linear independence of the solutions. Putting both conclusions together, it is therefore enough to test W
x at just one point of I to determine linear dependence or independence of these solutions. This gives us great latitude to choose a point at which the Wronskian is easy to evaluate.

EXAMPLE 2.2

In Example 2.1, we considered the solutions y1
x = cos
x and y2
x = sin
x of y

+ y = 0, for all x. In this case, linear independence was obvious. The Wronskian of these solutions is    cos
x sin
x    W
x =    − sin
x cos
x = cos2
x + sin2
x = 1 = 0

EXAMPLE 2.3

It is not always obvious whether two solutions are linearly independent or dependent on an interval. Consider the equation y

+ xy = 0. This equation appears simple but is not easy to solve. By a power series method we will develop later, we can write two solutions 1 1 6 1 y1
x = 1 − x3 + x − x9 + · · · 6 180 12,960 and y2
x = x −

1 4 1 7 1 x + x − x10 + · · ·  12 504 45,360

2.2 Theory of Solutions of y

+ p
xy
+ q
xy = f
x

67

with both series converging for all x. Here I is the entire real line. The Wronskian of these solutions at any nonzero x would be difficult to evaluate, but at x = 0 we easily obtain W
0 = y1
0y2

0 − y1

0y2
0 =
1
1 −
0
0 = 1 Nonvanishing of the Wronskian at this one point is enough to conclude linear independence of these solutions. We are now ready to use the machinery we have built up to determine what is needed to find all solutions of y

+ p
xy
+ q
x = 0. THEOREM 2.4

Let y1 and y2 be linearly independent solutions of y

+ p
xy
+ q
xy = 0 on an open interval I. Then, every solution of this differential equation on I is a linear combination of y1 and y2 . This fundamental theorem provides a strategy for finding all solutions of y

+ p
xy
+ q
xy = 0 on I. Find two linearly independent solutions. Depending on p and q, this may be difficult, but at least we have a specific goal. If necessary, use the Wronskian to test for independence. The general linear combination c1 y1 + c2 y2 , with c1 and c2 arbitrary constants, then contains all possible solutions. We will prove the theorem following introduction of some standard terminology.

DEFINITION 2.2

Let y1 and y2 be solutions of y

+ p
xy
+ q
xy = 0 on an open interval I. 1. y1 and y2 form a fundamental set of solutions on I if y1 and y2 are linearly independent on I. 2. When y1 and y2 form a fundamental set of solutions, we call c1 y1 + c2 y2 , with c1 and c2 arbitrary constants, the general solution of the differential equation on I.

In these terms, we find the general solution by finding a fundamental set of solutions. Here is a proof of Theorem 2.4. Let  be any solution of y

+ p
xy
+ q
xy = 0 on I. We want to show that there must be numbers c1 and c2 such that Proof


x = c1 y1
x + c2 y2
x Choose any x0 in I. Let 
x0  = A and 

x0  = B. By Theorem 2.1,  is the unique solution on I of the initial value problem y

+ p
xy
+ q
xy = 0

y
x0  = A y

x0  = B

Now consider the system of two algebraic equations in two unknowns: y1
x0 c1 + y2
x0 c2 = A y1

x0 c1 + y2

x0 c2 = B

68

CHAPTER 2

Second-Order Differential Equations

It is routine to solve these algebraic equations. Assuming that W
x0  = 0, we find that c1 =

Ay2

x0  − By2
x0  By
x  − Ay1

x0   c2 = 1 0  W
x0  W
x0 

With this choice of c1 and c2 , the function c1 y1 + c2 y2 is a solution of the initial value problem. By uniqueness of the solution of this problem, 
x = c1 y1
x + c2 y2
x on I, and the proof is complete. The proof reinforces the importance of having a fundamental set of solutions, since the nonvanishing of the Wronskian plays a vital role in showing that an arbitrary solution must be a linear combination of the fundamental solutions.

2.2.2

The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

The ideas just developed for the homogeneous equation (2.2) also provide the key to solving the nonhomogeneous equation y

+ p
xy
+ q
xy = f
x

(2.3)

THEOREM 2.5

Let y1 and y2 be a fundamental set of solutions of y

+ p
xy
+ q
xy = 0 on an open interval I. Let yp be any solution of equation (2.3) on I. Then, for any solution  of equation (2.3), there exist numbers c1 and c2 such that  = c1 y1 + c2 y2 + yp  This conclusion leads us to call c1 y1 + c2 y2 + yp the general solution of equation (2.3) and suggests the following strategy. To solve y

+ p
xy
+ q
xy = f
x: 1. find the general solution c1 y1 + c2 y2 of the associated homogeneous equation y

+ p
xy
+ q
xy = 0, 2. find any solution yp of y

+ p
xy
+ q
xy = f
x, and 3. write the general solution c1 y1 + c2 y2 + yp . This expression contains all possible solutions of equation (2.3) on the interval. Again, depending on p, q, and f , the first two steps may be formidable. Nevertheless, the theorem tells us what to look for and provides a clear way to proceed. Here is a proof of the theorem. Proof

Since  and yp are both solutions of equation (2.3), then
 − yp 

+ p
 − yp 
+ q
 − yp  = 

+ p
+ q −
yp

+ pyp
+ qyp  = f − f = 0

Therefore,  − yp is a solution of y

+ py
+ qy = 0. Since y1 and y2 form a fundamental set of solutions for this homogeneous equation, there are constants c1 and c2 such that  − yp = c1 y1 + c2 y2  and this is what we wanted to show.

2.3 Reduction of Order

69

The remainder of this chapter is devoted to techniques for carrying out the strategies just developed. For the general solution of the homogeneous equation (2.2) we must produce a fundamental set of solutions. And for the nonhomogeneous equation (2.3) we need to find one particular solution, together with a fundamental set of solutions of the associated homogeneous equation (2.2).

SECTION 2.2

PROBLEMS

In each of Problems 1 through 6, (a) verify that y1 and y2 are solutions of the differential equation, (b) show that their Wronskian is not zero, (c) write the general solution of the differential equation, and (d) find the solution of the initial value problem. 1. y

− 4y = 0 y
0 = 1 y

0 = 0 y1
x = cosh
2x y2
x = sinh
2x

9. Give an example to show that the product of two solutions of y

+ p
xy
+ q
xy = 0 need not be a solution. 10. Show that y1
x = 3e2x − 1 and y2
x = e−x + 2 are solutions of yy

+ 2y
−
y
2 = 0, but that neither 2y1 nor y1 + y2 is a solution. Why does this not contradict Theorem 2.2?

2. y

+ 9y = 0 y
/3 = 0 y

/3 = 1 y1
x = cos
3x y2
x = sin
3x 3. y

+ 11y
+ 24y = 0 y
0 = 1 y

0 = 4 y1
x = e−3x  y2
x = e−8x
4. y

+ 2y
+ 8y = 0 √ y
0 = 2 y
0 = 0 √ y1
x = e−x cos
7x y2
x = e−x sin
7x 7 16 5. y

− y
+ 2 y = 0 y
1 = 2 y

1 = 4 x x y1
x = x4  y2
x = x4 ln
x   1 1 6. y

+ y
+ 1 − 2 y = 0 y
 = −5 y

 = 8 x  4x  2 2 y1
x = cos
x y2
x = sin
x x x 7. Let y1
x = x2 and y2
x = x3 . Show that W
x = x4 for all real x. Then W
0 = 0, but W
x is not identically zero. Why does this not contradict Theorem 2.3.1, with the interval I chosen as the entire real line?

2.3

8. Show that y1
x = x and y2
x = x2 are linearly independent solutions of x2 y

− 2xy
+ 2y = 0 on −1 1, but that W
0 = 0. Why does this not contradict Theorem 2.3.1 on this interval?

11. Suppose y1 and y2 are solutions of y

+ p
xy
+ q
xy = 0 on a b, and that p and q are continuous on this interval. Suppose y1 and y2 both have a relative extremum at x0 in
a b. Prove that y1 and y2 are linearly dependent on a b. 12. Let  be a solution of y

+ p
xy
+ q
xy = 0 on an open interval I, and suppose 
x0  = 0 for some x0 in I. Suppose 
x is not identically zero. Prove that 

x0  = 0. 13. Let y1 and y2 be distinct solutions of y

+ p
xy
+ q
xy = 0 on an open interval I. Let x0 be in I and suppose y1
x0  = y2
x0  = 0. Prove that y1 and y2 are linearly dependent on I. Thus linearly independent solutions cannot share a common zero.

Reduction of Order Given y

+ p
xy
+ q
xy = 0, we want two independent solutions. Reduction of order is a technique for finding a second solution, if we can somehow produce a first solution. Suppose we know a solution y1 , which is not identically zero. We will look for a second solution of the form y2
x = u
xy1
x. Compute y2
= u
y1 + uy1


y2

= u

y1 + 2u
y1
+ uy1



In order for y2 to be a solution we need u

y1 + 2u
y1
+ uy1

+ p u
y1 + uy1
 + quy1 = 0

70

CHAPTER 2

Second-Order Differential Equations

Rearrange terms to write this equation as u

y1 + u
2y1
+ py1  + u y1

+ py1
+ qy1  = 0 The coefficient of u is zero because y1 is a solution. Thus we need to choose u so that u

y1 + u
2y1
+ py1  = 0 On any interval in which y1
x = 0, we can write u

+

2y1
+ py1
u = 0 y1

To help focus on the problem of determining u, denote g
x =

2y1

x + p
xy1
x  y1
x

a known function because y1
x and p
x are known. Then u

+ g
xu
= 0 Let v = u
to get v
+ g
xv = 0 This is a linear first-order differential equation for v, with general solution v
x = Ce−



g
x dx



Since we need only one second solution y2 , we will take C = 1, so v
x = e− Finally, since v = u
, u
x =




e−





g
x dx



g
x dx

dx

If we can perform these integrations and obtain u
x, then y2 = uy1 is a second solution of y

+ py
+ qy = 0. Further, W
x = y1 y2
− y1
y2 = y1
uy1
+ u
y1  − y1
uy1 = u
y12 = vy12  Since v
x is an exponential function, v
x = 0. And the preceding derivation was carried out on an interval in which y1
x = 0. Thus W
x = 0 and y1 and y2 form a fundamental set of solutions on this interval. The general solution of y

+ py
+ qy = 0 is c1 y1 + c2 y2 . We do not recommend memorizing formulas for g, v and then u. Given one solution y1 , substitute y2 = uy1 into the differential equation and, after the cancellations that occur because y1 is one solution, solve the resulting equation for u
x.

EXAMPLE 2.4

Suppose we are given that y1
x = e−2x is one solution of y

+ 4y
+ 4y = 0. To find a second solution, let y2
x = u
xe−2x . Then y2
= u
e−2x − 2e−2x u

and

y2

= u

e−2x + 4e−2x u − 4u
e−2x 

Substitute these into the differential equation to get u

e−2x + 4e−2x u − 4u
e−2x + 4
u
e−2x − 2e−2x u + 4ue−2x = 0

2.3 Reduction of Order

71

Some cancellations occur because e−2x is one solution, leaving u

e−2x = 0 or u

= 0 Two integrations yield u
x = cx + d. Since we only need one second solution y2 , we only need one u, so we will choose c = 1 and d = 0. This gives u
x = x and y2
x = xe−2x  Now

   e−2x xe−2x   W
x =   = e−4x = 0 −2e−2x e−2x − 2xe−2x 

for all x. Therefore, y1 and y2 form a fundamental set of solutions for all x, and the general solution of y

+ 4y
+ 4y = 0 is y
x = c1 e−2x + c2 xe−2x 

EXAMPLE 2.5

Suppose we want the general solution of y

−
3/xy
+
4/x2 y = 0 for x > 0, and somehow we find one solution y1
x = x2 . Put y2
x = x2 u
x and compute y2
= 2xu + x2 u


and

y2

= 2u + 4xu
+ x2 u



Substitute into the differential equation to get 3 4 2u + 4xu
+ x2 u

−
2xu + x2 u
 + 2
x2 u = 0 x x Then x2 u

+ xu
= 0 Since the interval of interest is x > 0, we can write this as xu

+ u
= 0 With v = u
, this is xv
+ v =
xv
= 0 so xv = c We will choose c = 1. Then v = u
=

1 x

so u = ln
x + d and we choose d = 0 because we need only one suitable u. Then y2
x = x2 ln
x is a second solution. Further, for x > 0,   2 x2 ln
x x = x3 = 0 W
x = 2x 2x ln
x + x

CHAPTER 2

72

Second-Order Differential Equations

Then x2 and x2 ln
x form a fundamental set of solutions for x > 0. The general solution is for x > 0 is y
x = c1 x2 + c2 x2 ln
x

PROBLEMS

SECTION 2.3

In each of Problems 1 through 10, verify that the given function is a solution of the differential equation, find a second solution by reduction of order, and finally write the general solution. 1. y + 4y = 0 y1
x = cos
2x 2. y − 9y = 0 y1
x = e






(d) y

+
y
2 = 0

3x

3. y − 10y + 25y = 0 y1
x = e

(e) y

= 1 +
y
2 5x

4. x2 y

− 7xy
+ 16y = 0 y1
x = x4 for x > 0 5. x2 y

− 3xy
+ 4y = 0 y1
x = x2 for x > 0 6.
2x2 + 1y

− 4xy
+ 4y = 0 y1
x = x for x > 0 1 8 7. y − y
− 2 y = 0 y1
x = x4 for x > 0 x x



8. y

−

(b) xy

+ 2y
= x (c) 1 − y
= 4y











(a) xy

= 2 + y


2x
2 y + y = 0 y1
x = x 1 + x2 1 + x2

  1
1 1 9. y + y + 1 − 2 y = 0 y1
x = √ cos
x for x 4x x x>0



10.
2x2 + 3x + 1y

+ 2xy
− 2y = 0 y1
x = x on any 1 interval not containing −1 or − 2 11. Verify that, for any nonzero constant a, y1
x = e−ax is a solution of y

+ 2ay
+ a2 y = 0. Write the general solution. 12. A second-order equation F
x y y
 y

 = 0 in which y is not explicitly present can sometimes be solved by putting u = y
. This results in a first-order equation G
x u u
 = 0. If this can be solved for u
x, then y1
x = u
x dx is a solution of the given secondorder equation. Use this method to find one solution, then find a second solution, and finally the general solution of the following.

13. A second-order equation in which x does not explicitly appear can sometimes be solved by putting u = y
and thinking of y as the independent variable and u as a function of y. Write du du dy du d dy

= = =u y = dx dx dx dy dx dy to convert F
y y
 y

 = 0 into the first-order equation F
y u u
du/dy = 0. Solve this equation for u
y and then set u = y
to solve for y as a function of x. Use this method to find a solution (perhaps implicitly defined) of each of the following. (a) yy

+ 3
y
2 = 0 (b) yy

+
y + 1
y
2 = 0 (c) yy

= y2 y
+
y
2 (d) y

= 1 +
y
2 (e) y

+
y
2 = 0 14. Consider y

+ Ay
+ By = 0, in which A and B are constants and A2 − 4B = 0. Show that y1
x = e−Ax/2 is one solution, and use reduction of order to find the second solution y2
x = xe−Ax/2 . 15. Consider y

+
A/xy
+
B/x2 y = 0 for x > 0, with A and B constants such that
A − 12 − 4B = 0. Verify that y1
x = x
1−A/2 is one solution, and use reduction of order to derive the second solution y2
x = x
1−A/2 ln
x.

2.4 The Constant Coefficient Homogeneous Linear Equation

2.4

73

The Constant Coefficient Homogeneous Linear Equation The linear homogeneous equation y

+ Ay
+ By = 0

(2.4)

in which A and B are numbers, occurs frequently in important applications. There is a standard approach to solving this equation. The form of equation (2.4) requires that constant multiples of derivatives of y
x must sum to zero. Since the derivative of an exponential function ex is a constant multiple of ex , we will look for solutions y
x = ex . To see how to choose , substitute ex into equation (2.4) to get 2 ex + Aex + Bex = 0 This can only be true if 2 + A + B = 0 This is called the characteristic equation of equation (2.4). Its roots are √ −A ± A2 − 4B =  2 leading to three cases.

2.4.1

Case 1: A2 − 4B > 0

In this case the characteristic equation has two real, distinct roots, √ √ −A − A2 − 4B −A + A2 − 4B a= and b = 2 2 yielding solutions y1
x = eax and y2
x = ebx for equation (2.4). These form a fundamental set of solutions on the real line, since W
x = eax bebx − ebx aebx =
b − ae
a+bx and this is nonzero because a = b. The general solution in this case is y
x = c1 eax + c2 ebx 

EXAMPLE 2.6

The characteristic equation of y

− y
− 6y = 0 is 2 −  − 6 = 0 with roots a = −2 and b = 3. The general solution is y = c1 e−2x + c2 e3x 

74

CHAPTER 2

2.4.2

Second-Order Differential Equations

Case 2: A2 − 4B = 0

Now the characteristic equation has the repeated root  = −A/2, so y1
x = e−Ax/2 is one solution. This method does not provide a second solution, but we have reduction of order for just such a circumstance. Try y2
x = u
xe−Ax/2 and substitute into the differential equation to get   A A2 −Ax/2 ue − Au
e−Ax/2 + u

e−Ax/2 + A − ue−Ax/2 + u
e−Ax/2 + Bue−Ax/2 = 0 4 2 Divide by e−Ax/2 and rearrange terms to get   A2

u + B− u = 0 4 Because in the current case, A2 − 4B = 0, this differential equation reduces to just u

= 0, and we can choose u
x = x. A second solution in this case is y2
x = xe−Ax/2 . Since y1 and y2 are linearly independent, they form a fundamental set and the general solution is y
x = c1 e−Ax/2 + c2 xe−Ax/2 = e−Ax/2
c1 + c2 x EXAMPLE 2.7

The characteristic equation of y

− 6y
+ 9y = 0 is 2 − 6 + 9 = 0, with repeated root  = 3. The general solution is y
x = e3x
c1 + c2 x

2.4.3

Case 3: A2 − 4B < 0

Now the characteristic equation has complex roots √ −A ± 4B − A2 i  2 For convenience, write A p=−  2

q=

1√ 4B − A2  2

so the roots of the characteristic equation are p ± iq. This yields two solutions y1
x = e
p+iqx

and y2
x = e
p−iqx 

These are linearly independent because their Wronskian is    e
p+iqx e
p−iqx   W
x = 
p + iqe
p+iqx
p − iqe
p−iqx  =
p − iqe2px −
p + iqe2px = −2iqe2px  and this is nonzero in the current case in which q = 0. Therefore the general solution is y
x = c1 e
p+iqx + c2 e
p−iqx 

(2.5)

EXAMPLE 2.8

√ The characteristic equation of y

+ 2y
+ 6y = 0 is 2 + 2 + 6 = 0, with roots −1 ± 5i. The general solution is y
x = c1 e
−1+

√ 5ix

+ c2 e
−1−

√ 5ix



2.4 The Constant Coefficient Homogeneous Linear Equation

2.4.4

75

An Alternative General Solution in the Complex Root Case

When the characteristic equation has complex roots, we can write a general solution in terms of complex exponential functions. This is sometimes inconvenient, for example, in graphing the solutions. But recall that any two linearly independent solutions form a fundamental set. We will therefore show how to use the general solution (2.5) to find a fundamental set of real-valued solutions. Begin by recalling the Maclaurin expansions of ex , cos
x, and sin
x: ex =

 1 n 1 1 1 1 x = 1 + x + x2 + x3 + x4 + x5 + · · ·  n! 2! 3! 4! 5! n=0

cos
x =


−1n n=0


2n!

x2n = 1 −

1 2 1 4 1 6 x + x − x +···  2! 4! 6!

and sin
x =


−1n 2n+1 1 1 1 x = x − x3 + x5 − x7 + · · ·  3! 5! 7! n=0
2n + 1!

with each series convergent for all real x. The eighteenth century Swiss mathematician Leonhard Euler experimented with replacing x with ix in the exponential series and noticed an interesting relationship between the series for ex , cos
x, and sin
x. First, eix =

 1
ixn n! n=0

= 1 + ix +

1 1 1 1 1
ix2 +
ix3 +
ix4 +
ix5 +
ix6 + · · ·  2! 3! 4! 5! 6!

Now, integer powers of i repeat the values i, −1, −i, 1 with a period of four: i2 = −1

i3 = −i

i4 = 1

i5 = i4 i = i

i6 = i4 i2 = −1

i7 = i4 i3 = −i

and so on, continuing in cyclic fashion. Using this fact in the Maclaurin series for eix , we obtain 1 i 1 i 1 eix = 1 + ix − x2 − x3 + x4 + x5 − x6 − · · · 2! 3! 4! 5! 6!     1 3 1 5 1 2 1 4 1 6 = 1− x + x − x +··· +i x− x + x −··· 2! 4! 6! 3! 5! = cos
x + i sin
x

(2.6)

This is Euler’s formula. In a different form, it was discovered a few years earlier by Newton’s contemporary Roger Cotes (1682–1716). Cotes is not of the stature of Euler, but Newton’s high opinion of him is reflected in Newton’s remark, “If Cotes had lived, we would have known something.” Since cos
−x = cos
x and sin
−x = − sin
x, replacing x by −x in Euler’s formula yields e−ix = cos
x − i sin
x Now return to the problem of solving y

+ Ay
+ By = 0 when the characteristic equation has complex roots p ± iq. Since p and q are real numbers, we have e
p+iqx = epx eiqx = epx
cos
qx + i sin
qx = epx cos
qx + iepx sin
qx

76

CHAPTER 2

Second-Order Differential Equations

and e
p−iqx = epx e−iqx = epx
cos
qx − i sin
qx = epx cos
qx − iepx sin
qx The general solution (2.5) can therefore be written y
x = c1 epx cos
qx + ic1 epx sin
qx + c2 epx cos
qx − ic2 epx sin
qx =
c1 + c2 epx cos
qx +
c1 − c2 iepx sin
qx We obtain solutions for any numerical choices of c1 and c2 . In particular, if we choose 1 c1 = c2 = we obtain the solution 2 y3
x = epx cos
qx And if we put c1 =

1 1 and c2 = − we obtain still another solution 2i 2i y4
x = epx sin
qx

Further, these last two solutions are linearly independent on the real line, since    epx cos
qx epx sin
qx  W
x =   pepx cos
qx − qepx sin
qx pepx sin
qx + qepx cos
qx = e2px
sin
qx cos
qx + cos2
qx − sin
qx cos
qx + sin2
qx = e2px = 0

for all real x

We can therefore, if we prefer, form a fundamental set of solutions using y3 and y4 , writing the general solution of y

+ Ay
+ By = 0 in this case as y
x = epx
c1 cos
qx + c2 sin
qx This is simply another way of writing the general solution of equation (2.4) in the complex root case.

EXAMPLE 2.9

Revisiting the equation y

− 6y
+ 6y = 0 of Example 2.8, we can also write the general solution  √ √ y
x = e−x c1 cos 5x + c2 sin 5x  We now have the general solution of the constant coefficient linear homogeneous equation y

+ Ay
+ By = 0 in all cases. As usual, we can solve an initial value problem by first finding the general solution of the differential equation, then solving for the constants to satisfy the initial conditions.

EXAMPLE 2.10

Solve the initial value problem y

− 4y
+ 53y = 0

y
 = −3 y

 = 2

First solve the differential equation. The characteristic equation is 2 − 4 + 53 = 0

2.4 The Constant Coefficient Homogeneous Linear Equation

77

with complex roots 2 ± 7i. The general solution is y
x = c1 e2x cos
7x + c2 e2x sin
7x Now y
 = c1 e2 cos
7 + c2 e2 sin
7 = −c1 e2 = −3 so c1 = 3e−2  Thus far y
x = 3e−2 e2x cos
7x + c2 e2x sin
7x Compute y

x = 3e−2 2e2x cos
7x − 7e2x sin
7x + 2c2 e2x sin
7x + 7c2 e2x cos
7x Then y

 = 3e−2 2e2
−1 + 7c2 e2
−1 = 2 so 8 c2 = − e−2  7 The solution of the initial value problem is 8 y
x = 3e−2 e2x cos
7x − e−2 e2x sin
7x 7

8 2
x− 3 cos
7x − sin
7x  =e 7

SECTION 2.4

PROBLEMS

In each of Problems 1 through 12, find the general solution.

In each of Problems 13 through 21, solve the initial value problem.

1. y

− y
− 6y = 0

13. y

+ 3y
= 0 y
0 = 3 y

0 = 6

2. y

− 2y
+ 10y = 0

14. y

+ 2y
− 3y = 0 y
0 = 6 y

0 = −2

3. y

+ 6y
+ 9y = 0

15. y

− 2y
+ y = 0 y
1 = y

1 = 0

4. y

− 3y
= 0

16. y

− 4y
+ 4y = 0 y
0 = 3 y

0 = 5








5. y + 10y + 26y = 0






6. y + 6y − 40y = 0 7. y

+ 3y
+ 18y = 0 8. y

+ 16y
+ 64y = 0 9. y

− 14y
+ 49y = 0 10. y

− 6y
+ 7y = 0 11. y

+ 4y
+ 9y = 0 12. y

+ 5y
= 0

17. y

+ y
− 12y = 0 y
2 = 2 y

2 = 1 18. y

− 2y
− 5y = 0 y
0 = 0 y

0 = 3 19. y

− 2y
+ y = 0 y
1 = 12 y

1 = −5 20. y

− 5y
+ 12y = 0 y
2 = 0 y

2 = −4 21. y

− y
+ 4y = 0 y
−2 = 1 y

−2 = 3 22. This problem illustrates how small changes in the coefficients of a differential equation may cause dramatic changes in the solutions.

CHAPTER 2

78

Second-Order Differential Equations

(a) Find the general solution 
x of y

− 2ay
+ a2 y = 0, with a a nonzero constant.

(b) Find the solution  of the initial value problem

(b) Find the general solution 
x of y

− 2ay
+
a2 − 2 y = 0, in which  is a positive constant.

Here  is any positive number.

y

− 2ay
+
a2 − 2 y = 0

(c) Is it true that lim→0 
x = 
x? How does this answer differ, if at all, from the conclusion in Problem 22(c)?

(c) Show that, as  → 0, the differential equation in (b) approaches in a limit sense the differential equation in (a), but the solution 
x for (b) does not in general approach the solution 
x for (a).

24. Suppose  is a solution of

23. (a) Find the solution  of the initial value problem






y − 2ay + a y = 0 2

y

+ Ay
+ By = 0




y
0 = c y
0 = d

y
x0  = a y

x0  = b

Here A, B, a, and b are constants. Suppose A and B are positive. Prove that limx→ 
x = 0.

with a, c, and d constants and a = 0.

2.5

y
0 = c y

0 = d

Euler’s Equation In this section we will define another class of second-order differential equations for which there is an elementary technique for finding the general solution. The second-order homogeneous equation 1 1 y

+ Ay
+ 2 By = 0 x x

(2.7)

with A and B constant, is called Euler’s equation. It is defined on the half-lines x > 0 and x < 0. We will assume for this section that x > 0. We will solve Euler’s equation by transforming it to a constant coefficient linear equation, which we can solve easily. Recall that any positive number x can be written as et for some t (namely for t = ln
x). Make the change of variables x = et 

or, equivalently, t = ln
x

and let Y
t = y
et  That is, in the function y
x, replace x by et , obtaining a new function of t. For example, if y
x = x3 , then Y
t =
et 3 = e3t . Now compute chain-rule derivatives. First, y

x =

dY dt 1 = Y

t dt dx x

so Y

t = xy

x

2.5 Euler’s Equation Next, d
d y
x = y
x = dx dx





79



1
Y
t x

1 d
1
Y
t Y
t + 2 x x dx 1 1 dY
dt = − 2 Y

t + x x dt dx 1 1 1
= − 2 Y
t + Y


t x x x 1

= 2
Y
t − Y

t x =−

Therefore, x2 y


x = Y


t − Y

t If we write Euler’s equation as x2 y


x + Axy

x + By
x = 0 then these substitutions yield Y


t − Y

t + AY

t + BY
t = 0 or Y

+
A − 1Y
+ BY = 0

(2.8)

This is a constant coefficient homogeneous linear differential equation for Y
t. Solve this equation, then let t = ln
x in the solution Y
t to obtain y
x satisfying the Euler equation. We need not repeat this derivation each time we want to solve an Euler equation, since the coefficients A − 1 and B for the transformed equation (2.8) are easily read from the Euler equation (2.7). In carrying out this strategy, it is useful to recall that, for x > 0, xr = er ln
x 

EXAMPLE 2.11

Find the general solution of x2 y

+ 2xy
− 6y = 0. Upon letting x = et , this differential equation transforms to Y

+ Y
− 6Y = 0 The coefficient of Y
is A − 1, with A = 2 in Euler’s equation. The general solution of this linear homogeneous differential equation is Y
t = c1 e−3t + c2 e2t for all real t. Putting t = ln
x with x > 0, we obtain y
x = c1 e−3 ln
x + c2 e2 ln
x = c1 x−3 + c2 x2  and this is the general solution of the Euler equation.

80

CHAPTER 2

Second-Order Differential Equations

EXAMPLE 2.12

Consider the Euler equation x2 y

− 5xy
+ 9y = 0. The transformed equation is Y

− 6Y
+ 9Y = 0 with general solution Y
t = c1 e3t + c2 te3t  Let t = ln
x to obtain y
x = c1 x3 + c2 x3 ln
x for x > 0. This is the general solution of the Euler equation.

EXAMPLE 2.13

Solve x2 y

+ 3xy
+ 10y = 0 This transforms to Y

+ 2Y
+ 10Y = 0 with general solution Y
t = c1 e−t cos
3t + c2 e−t sin
3t Then y
x = c1 x−1 cos
3 ln
x + c2 x−1 sin
3 ln
x 1 =
c1 cos
3 ln
x + c2 sin
3 ln
x x for x > 0. As usual, we can solve an initial value problem by finding the general solution of the differential equation, then solving for the constants to satisfy the initial conditions.

EXAMPLE 2.14

Solve the initial value problem x2 y

− 5xy
+ 10y = 0

y
1 = 4 y

1 = −6

We will first find the general solution of the Euler equation, then determine the constants to satisfy the initial conditions. With t = ln
x, we obtain Y

− 6Y
+ 10Y = 0 having general solution Y
t = c1 e3t cos
t + c2 e3t sin
t The general solution of the Euler equation is y
x = c1 x3 cos
ln
x + c2 x3 sin
ln
x

2.5 Euler’s Equation

81

For the first initial condition, we need y
1 = 4 = c1  Thus far, y
x = 4x3 cos
ln
x + c2 x3 sin
ln
x Then y

x = 12x2 cos
ln
x − 4x2 sin
ln
x + 3c2 x2 sin
ln
x + c2 x2 cos
ln
x so y

1 = 12 + c2 = −6 Then c2 = −18 and the solution of the initial value problem is y
x = 4x3 cos
ln
x − 18x3 sin
ln
x

Observe the structure of the solutions of different kinds of differential equations. Solutions of the constant coefficient linear equation y

+ Ay
+ By = 0 must have the forms ex , xex , ex cos
x, or ex sin
x, depending on the coefficients. And solutions of an Euler equation x2 y

+ Axy
+ By = 0 must have the forms xr , xr ln
x, xp cos
q ln
x, or xp sin
q ln
x. For example, x3 could never be a solution of the linear equation and e−6x could never be the solution of an Euler equation.

SECTION 2.5

PROBLEMS

In each of Problems 1 through 12, find the general solution. 1. x2 y

+ 2xy
− 6y = 0 2. x2 y

+ 3xy
+ y = 0 3. x2 y

+ xy
+ 4y = 0 4. x2 y

+ xy
− 4y = 0 5. x2 y

+ xy
− 16y = 0 6. x2 y

+ 3xy
+ 10y = 0 7. x2 y

+ 6xy
+ 6y = 0

13. x2 y

+ 5xy
+ 20y = 0 y
−1 = 3 y

−1 = 2 (Here the solution of Euler’s equation for x < 0 is needed). 14. x2 y

+ 5xy
− 21y = 0 y
2 = 1 y

2 = 0 15. x2 y

− xy
= 0 y
2 = 5 y

2 = 8 16. x2 y

− 3xy
+ 4y = 0 y
1 = 4 y

1 = 5 17. x2 y

+ 7xy
+ 13y = 0 y
−1 = 1 y

−1 = 3 18. x2 y

+ xy
− y = 0 y
2 = 1 y

2 = −3

8. x2 y

− 5xy
+ 58y = 0

19. x2 y

+ 25xy
+ 144y = 0 y
1 = −4 y

1 = 0

9. x2 y

+ 25xy
+ 144y = 0

20. x2 y

− 9xy
+ 24y = 0 y
1 = 1 y

1 = 10

10. x2 y

− 11xy
+ 35y = 0

21. x2 y

+ xy
− 4y = 0 y
1 = 7 y

1 = −3

11. x2 y

− 2xy
+ 12y = 0

22. Here is another approach to solving an Euler equation. For x > 0, substitute y = xr and obtain values of r to make this a solution. Show how this leads in all cases to the same general solution as obtained by the transformation method.

2



12. x y + 4y = 0 In each of Problems 13 through 21, solve the initial value problem.

82

2.6

CHAPTER 2

Second-Order Differential Equations

The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x In view of Theorem 2.5, if we are able to find the general solution yh of the linear homogeneous equation y

+ p
xy
+ q
xy = 0, then the general solution of the linear nonhomogeneous equation y

+ p
xy
+ q
xy = f
x

(2.9)

is y = yh + yp , in which yp is any solution of equation (2.9). This section is devoted to two methods for finding such a particular solution yp .

2.6.1

The Method of Variation of Parameters

Suppose we can find a fundamental set of solutions y1 and y2 for the homogeneous equation. The general solution of this homogeneous equation has the form yh
x = c1 y1
x + c2 y2
x. The method of variation of parameters consists of attempting a particular solution of the nonhomogeneous equation by replacing the constants c1 and c2 with functions of x. Thus, attempt to find u
x and v
x so that yp
x = u
xy1
x + v
xy2
x is a solution of equation (2.9). How should we choose u and v? First compute yp
= uy1
+ vy2
+ u
y1 + v
y2  In order to simplify this expression, the first condition we will impose on u and v is that u
y1 + v
y2 = 0

(2.10)

Now yp
= uy1
+ vy2
 Next compute yp

= u
y1
+ v
y2
+ uy1

+ vy2

 Substitute these expressions for yp
and yp

into equation (2.9): u
y1
+ v
y2
+ uy1

+ vy2

+ p
x
uy1
+ vy2
 + q
x
uy1 + vy2  = f
x Rearrange terms in this equation to get u y1

+ p
xy1
+ q
xy1  + v y2

+ p
xy2
+ q
xy2  + u
y1
+ v
y2
= f
x The two terms in square brackets vanish because y1 and y2 are solutions of the homogeneous equation. This leaves u
y1
+ v
y2
= f
x

(2.11)

Now solve equations (2.10) and (2.11) for u
and v
to get u

x = −

y2
xf
x W
x

and v

x =

y1
xf
x W
x

(2.12)

2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

83

in which W is the Wronskian of y1 and y2 . If we can integrate these equations to determine u and v, then we have yp .

EXAMPLE 2.15

We will find the general solution of y

+ y = sec
x for −/4 < x < /4. The characteristic equation of y

+ 4y = 0 is 2 + 4 = 0, with roots ±2i. We may therefore choose y1
x = cos
2x and y2
x = sin
2x. The Wronskian of these solutions of the homogeneous equation is    cos
2x sin
2x  W
x =   = 2 −2 sin
2x 2 cos
2x With f
x = sec
x, equations (2.12) give us 1 u

x = − sin
2x sec
x 2 1 1 = − 2 sin
x cos
x = − sin
x 2 cos
x and 1 1 1 v

x = cos
2x sec
x = 2 cos2
x − 1 2 2 cos
x = cos
x − Then u
x = and v
x =




1 sec
x 2


− sin
x dx = cos
x

cos
x dx −

= sin
x −

1
sec
x dx 2

1 ln sec
x + tan
x 2

Here we have let the constants of integration be zero because we need only one u and one v. Now we have the particular solution yp
x = u
xy1
x + v
xy2
x   1 = cos
x cos
2x + sin
x − ln sec
x + tan
x sin
2x 2 The general solution of y

+ y = sec
x is y
x = yh
x + yp
x = c1 cos
2x + c2 sin
2x   1 + cos
x cos
2x + sin
x − ln sec
x + tan
x sin
2x 2

84

CHAPTER 2

Second-Order Differential Equations

EXAMPLE 2.16

Suppose we want the general solution of 4 4 y

− y
+ 2 y = x2 + 1 x x for x > 0. The associated homogeneous equation is 4 4 y

− y
+ 2 y = 0 x x which we recognize as an Euler equation, with fundamental solutions y1
x = x and y2
x = x4 for x > 0. The Wronskian of these solutions is   x x4   = 3x4 W
x =  1 4x3  and this is nonzero for x > 0. From equations (2.12), u

x = −

x4
x2 + 1 1 = −
x2 + 1 4 3x 3

and x
x2 + 1 1 v
x = = 3x4 3




 1 1 +  x x3

Integrate to get 1 1 u
x = − x3 − x 9 3 and v
x =

1 1 ln
x − 2  3 6x

A particular solution is     1 1 1 1 ln
x − 2 x4  yp
x = − x3 − x x + 9 3 3 6x The general solution is y
x = yh
x + yp
x 1 1 1 1 = c1 x + c2 x4 − x4 − x2 + x4 ln
x − x2 9 3 3 6 1 1 1 = c1 x + c2 x4 − x4 − x2 + x4 ln
x 9 2 3 for x > 0

2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

2.6.2

85

The Method of Undetermined Coefficients

Here is a second method for finding a particular solution yp , but it only applies if p
x and q
x are constant. Thus consider y

+ Ay
+ By = f
x Sometimes we can guess the general form of a solution yp from the form of f
x. For example, suppose f
x is a polynomial. Since derivatives of polynomials are polynomials, we might try a polynomial for yp
x Substitute a polynomial with unknown coefficients into the differential equation, and then choose the coefficients to match y

+ Ay
+ By with f
x. Or suppose f
x is an exponential function, say f
x = e−2x . Since derivatives of e−2x are just constant multiples of e−2x , we would attempt a solution of the form yp = Ce−2x , substitute into the differential equation, and solve for C to match the left and right sides of the differential equation. Here are some examples of this method.

EXAMPLE 2.17

Solve y

− 4y = 8x2 − 2x. Since f
x = 8x2 − 2x is a polynomial of degree 2, we will attempt a solution yp
x = ax2 + bx + c We do not need to try a higher degree polynomial, since the degree of y

− 4y must be 2. If, for example, we included an x3 term in yp , then yp

− 4yp would have an x3 term, and we know that it does not. Compute yp
= 2ax + b

and yp

= 2a

and substitute into the differential equation to get 2a − 4
ax2 + bx + c = 8x2 − 2x Collect coefficients of like powers of x to write
−4a − 8x2 +
−4b + 2x +
2a − 4c = 0 For yp to be a solution for all x, the polynomial on the left must be zero for all x. But a second degree polynomial can have only two roots, unless it is the zero polynomial. Thus all the coefficients must vanish, and we have the equations −4a − 8 = 0 −4b + 2 = 0 2a − 4c = 0 Solve these to obtain a = −2

1 b=  2

c = −1

Thus a solution is 1 yp
x = −2x2 + x − 1 2 as can be verified by substitution into the differential equation.

86

CHAPTER 2

Second-Order Differential Equations

If we want the general solution of the differential equation, we need the general solution yh of y

− 4y = 0. This is yh
x = c1 e2x + c2 e−2x  The general solution of y

− 4y = 8x2 − 2x is 1 y
x = c1 e2x + c2 e−2x − 2x2 + x − 1 2 The method we have just illustrated is called the method of undetermined coefficients, because the idea is to guess a general form for yp and then solve for the coefficients to make a solution. Here are two more examples, after which we will point out a circumstance in which we must supplement the method.

EXAMPLE 2.18

Solve y

+ 2y
− 3y = 4e2x . Because f
x is a constant times an exponential, and the derivative of such a function is always a constant times the same function, we attempt yp = ae2x . Then yp
= 2ae2x and yp

= 4ae2x . Substitute into the differential equation to get 4ae2x + 4ae2x − 3ae2x = 4e2x  4 to get the solution 5 4 yp
x = e2x  5 Again, if we wish we can write the general solution

Then 5ae2x = 4e2x , so choose a =

4 y
x = c1 e−3x + c2 ex + e2x  5

EXAMPLE 2.19

Solve y

− 5y
+ 6y = −3 sin
2x. Here f
x = −3 sin
2x. Now we must be careful, because derivatives of sin
2x can be multiples sin
2x or cos
2x, depending on how many times we differentiate. This leads us to include both possibilities in a proposed solution: yp
x = c cos
2x + d sin
2x Compute yp
= −2c sin
2x + 2d cos
2x

yp

= −4c cos
2x − 4d sin
2x

Substitute into the differential equation to get − 4c cos
2x − 4d sin
2x − 5 −2c sin
2x + 2d cos
2x + 6 c cos
2x + d sin
2x = −3 sin
2x Collecting the cosine terms on one side and the sine terms on the other:

2d + 10c + 3 sin
2x = −2c + 10d cos
2x

2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

87

For yp to be a solution for all real x, this equation must hold for all x. But sin
2x and cos
2x are linearly independent (they are solutions of y

+ 4y = 0, and their Wronskian is nonzero). Therefore neither can be a constant multiple of the other. The only way the last equation can hold for all x is for the coefficient to be zero on both sides: 2d + 10c = −3 and 10d − 2c = 0 Then d=−

3 52

and c = −

15 52

and we have found a solution: yp
x = −

15 3 sin
2x − cos
2x 52 52

The general solution of this differential equation is y
x = c1 e3x + c2 e2x −

15 3 sin
2x − cos
2x 52 52

As effective as this method is, there is a difficulty that is intrinsic to the method. It can be successfully met, but one must be aware of it and know how to proceed. Consider the following example.

EXAMPLE 2.20

Solve y

+ 2y
− 3y = 8ex . The coefficients on the left side are constant, and f
x = 8ex seems simple enough, so we proceed with yp
x = cex . Substitute into the differential equation to get cex + 2cex − 3cex = 8ex  or 0 = 8ex  Something is wrong. What happened? The problem in this example is that ex is a solution of y

+ 2y
− 3y = 0, so if we substitute ce into y

+ 2y
− 3y = 8ex , the left side will equal zero, not 8ex . This difficulty will occur whenever the proposed yp contains a term that is a solution of the homogeneous equation y

+ Ay
+ By = 0, because then this term (which may be all of the proposed yp ) will vanish when substituted into y

+ Ay
+ By. There is a way out of this difficulty. If a term of the proposed yp is a solution of y

+ Ay
+ By = 0, multiply the proposed solution by x and try the modified function as yp . If this also contains a term, or by itself, satisfies y

+ Ay
+ By = 0, then multiply by x again to x

88

CHAPTER 2

Second-Order Differential Equations

try x2 times the original proposed solution. This is as far as we will have to go in the case of second-order differential equations. Now continue Example 2.20 with this strategy.

EXAMPLE 2.21

Consider again y

+ 2y
− 3y = 8ex . We saw that yp = cex does not work, because ex , and hence also cex , satisfies y

+ 2y
− 3y = 0. Try yp = cxex . Compute yp
= cex + cxex 

yp

= 2cex + cxex

and substitute into the differential equation to get 2cex + cxex + 2 cex + cxex  − 3cxex = 8ex  Some terms cancel and we are left with 4cex = 8ex  Choose c = 2 to obtain the particular solution yp
x = 2xex .

EXAMPLE 2.22

Solve y

− 6y
+ 9y = 5e3x . Our first impulse is to try yp = ce3x . But this is a solution of y

− 6y
+ 9y = 0. If we try yp = cxe3x , we also obtain an equation that cannot be solved for c. The reason is that the characteristic equation of y

− 6y
+ 9y = 0 is
 − 32 = 0, with repeated root 3. This means that e3x and xe3x are both solutions of the homogeneous equation y

− 6y
+ 9y = 0. Thus try yp
x = cx2 e3x . Compute yp
= 2cxe3x + 3cx2 e3x 

yp

= 2ce3x + 12cxe3x + 9cx2 e3x 

Substitute into the differential equation to get 2ce3x + 12cxe3x + 9cx2 e3x − 6 2cxe3x + 3cx2 e3x  + 9cx2 e3x = 5e3x  After cancellations we have 2ce3x = 5e3x so c = 5/2. We have found a particular solution yp
x = 5x2 e3x /2. The last two examples suggest that in applying undetermined coefficients to y

+Ay
+By = f
x, we should first obtain the general solution of y

+ Ay
+ By = 0. We need this anyway for a general solution of the nonhomogeneous equation, but it also tells us whether to multiply our first choice for yp by x or x2 before proceeding. Here is a summary of the method of undetermined coefficients. 1. From f
x, make a first conjecture for the form of yp . 2. Solve y

+ Ay
+ By = 0. If a solution of this equation appears in any term of the conjectured form for yp , modify this form by multiplying it by x. If this modified function still occurs in a solution of y

+ Ay
+ By = 0, multiply by x again (so the original yp is multiplied by x2 in this case). 3. Substitute the final proposed yp into y

+ Ay
+ By = f
x and solve for its coefficients.

2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

89

Here is a list of functions to try in the initial stage (1) of formulating yp . In this list P
x indicates a given polynomial of degree n, and Q
x and R
x polynomials with undetermined coefficients, of degree n. f
x

Initial Guess for yp

P
x

Q
x

ceax

deax

 cos
bx or  sin
bx

c cos
bx + d sin
bx

P
xeax

Q
xeax

P
x cos
bx or P
x sin
bx ax

ax

P
xe cos
bx or P
xe sin
bx

Q
x cos
bx + R
x sin
bx Q
xeax cos
bx + R
xeax sin
bx

EXAMPLE 2.23

Solve y

+ 9y = −4x sin
3x. With f
x = −4x sin
3x, the preceding list suggests that we attempt a particular solution of the form yp
x =
ax + b cos
3x +
cx + d sin
3x Now solve y

+ 9y = 0 to obtain the fundamental set of solutions cos
3x and sin
3x. The proposed yp includes terms b cos
3x and d sin
3x, which are also solutions of y

+ 9y = 0. Therefore, modify the proposed yp by multiplying it by x, trying instead yp
x =
ax2 + bx cos
3x +
cx2 + dx sin
3x Compute yp
=
2ax + b cos
3x −
3ax2 + 3bx sin
3x +
2cx + d sin
3x +
3cx2 + 3dx cos
3x and yp

=2a cos
3x −
6ax + 3b sin
3x −
6ax + 3b sin
3x −
9ax2 + 9bx cos
3x + 2c sin
3x +
6cx + 3d cos
3x +
6cx + 3d cos
3x −
9cx2 + 9dx sin
3x Substitute these into the differential equation to obtain 2a cos
3x −
6ax + 3b sin
3x −
6ax + 3b sin
3x −
9ax2 + 9bx cos
3x + 2c sin
3x +
6cx + 3d cos
3x +
6cx + 3d cos
3x −
9cx2 + 9dx sin
3x +
9ax2 + 9bx cos
3x +
9cx2 + 9dx sin
3x = −4x sin
3x Now collect coefficients of “like” terms
sin
3x, x sin
3x, x2 sin
3x, and so on). We get
2a + 6d cos
3x +
−6b + 2c sin
3x + 12cx cos
3x +
−12a + 4x sin
3x = 0

90

CHAPTER 2

Second-Order Differential Equations

with all other terms canceling. For this linear combination of cos
3x, sin
3x, x cos
3x, and x sin
3x to be zero for all x, each coefficient must be zero. Therefore, 2a + 6d = 0 −6b + 2c = 0 12c = 0 and −12a + 4 = 0 1 1 Then a = , c = 0, b = 0, and d = − . We have found the particular solution 3 9 1 1 yp
x = x2 cos
3x − x sin
3x 3 9 The general solution is 1 1 y
x = c1 cos
3x + c2 sin
3x + x2 cos
3x − x sin
3x 3 9 Sometimes a differential equation has nonconstant coefficients but transforms to a constant coefficient equation. We may then be able to use the method of undetermined coefficients on the transformed equation and then use the results to obtain solutions of the original equation.

EXAMPLE 2.24

Solve x2 y

− 5xy
+ 8y = 2 ln
x. The method of undetermined coefficients does not apply here, since the differential equation has nonconstant coefficients. However, from our experience with the Euler equation, apply the transformation t = ln
x and let Y
t = y
et . Using results from Section 2.5, the differential equation transforms to Y


t − 6Y

t + 8Y
t = 2t which has constant coefficients on the left side. The homogeneous equation Y

− 6Y
+ 8Y = 0 has general solution Yh
t = c1 e2t + c2 e4t and, by the method of undetermined coefficients, we find one solution of Y

− 6Y
+ 8Y = 2t to be 3 1 Yp
t = t +  4 16 The general solution for Y is 3 1 Y
t = c1 e2t + c2 e4t + t +  4 16 Since t = ln
x, the original differential equation for y has general solution 3 1 ln
x + 4 16 1 3 = c1 x2 + c2 x4 + ln
x +  4 16

y
x = c1 e2 ln
x + c2 e4 ln
x +

2.6 The Nonhomogeneous Equation y

+ p
xy
+ q
xy = f
x

2.6.3

91

The Principle of Superposition

Consider the equation y

+ p
xy
+ q
xy = f1
x + f2
x + · · · + fN
x

(2.13)

Suppose ypj is a solution of y

+ p
xy
+ q
xy = fj
x We claim that yp1 + yp2 + · · · + ypN is a solution of equation (2.13). This is easy to check by direct substitution into the differential equation:
yp1 + yp2 + · · · + ypN 

+ p
x
yp1 + yp2 + · · · + ypN 
+q
x
yp1 + yp2 + · · · + ypN 





=
yp1 + p
xyp1 + q
xyp1  + · · · +
ypN + p
xypN + q
xypN 

= f1
x + f2
x + · · · + fN
x This means that we can solve each equation y

+ p
xy
+ q
xy = fj
x individually, and the sum of these solutions is a solution of equation (2.13). This is called the principle of superposition, and it sometimes enables us to solve a problem by breaking it into a sum of “smaller” problems that are easier to handle individually.

EXAMPLE 2.25

Solve y

+ 4y = x + 2e−2x . Consider two problems: Problem 1: y

+ 4y = x, and Problem 2: y

+ 4y = 2e−2x . Using undetermined coefficients, we find that a solution of Problem 1 is yp1
x = x/4, and that a solution of Problem 2 is yp2
x = e−2x /4. Therefore, 1 yp
x =
x + e−2x  4 is a solution of y

+ 4y = x + 2e−2x . The general solution of this differential equation is 1 y
x = c1 cos
2x + c2 sin
2x +
x + e−2x  4

2.6.4

Higher-Order Differential Equations

The methods we now have for solving y

+ p
xy
+ q
xy = f
x under certain conditions can also be applied to higher-order differential equations, at least in theory. However, there are practical difficulties to this approach. Consider the following example.

92

CHAPTER 2

Second-Order Differential Equations

EXAMPLE 2.26

Solve d4 y dy d6 y + 15y = 0 −4 4 +2 6 dx dx dx If we take a cue from the second-order case, we attempt solutions y = ex . Upon substituting this into the differential equation, we obtain an equation for : 6 − 44 + 2 + 15 = 0 In the second-order case, the characteristic polynomial is always of degree 2 and easily solved. Here we encounter a sixth-degree polynomial whose roots are not obvious. They are, approximately, − 1685798616 ± 02107428331i − 004747911354 ± 1279046854i and 1733277730 ± 04099384482i

When the order of the differential equation is n > 2, having to find the roots of an nth degree polynomial is enough of a barrier to make this approach impractical except in special cases. A better approach is to convert this sixth-order equation to a system of first-order equations as follows. Define new variables z1 = y

z2 = y


z3 = y



z4 =

d3 y  dx3

z5 =

d4 y  dx4

z6 =

d5 y  dx5

Now we have a system of six first-order differential equations: z
1 = z2 z
2 = z3 z
3 = z4 z
4 = z5 z
5 = z6 z
6 = 4z5 − 2z2 − 15z1  The last equation in this system is exactly the original differential equation, stated in terms of the new quantities zj . The point to reformulating the problem in this way is that powerful matrix techniques can be invoked to find solutions. We therefore put off discussion of differential equations of order higher than 2 until we have developed the matrix machinery needed to exploit this approach.

2.7 Application of Second-Order Differential Equations to a Mechanical System

SECTION 2.6

PROBLEMS

In each of Problems 1 through 6, find the general solution using the method of variation of parameters. 1. y

+ y = tan
x 2. y

− 4y
+ 3y = 2 cos
x + 3 3. y

+ 9y = 12 sec
3x 4. y

− 2y
− 3y = 2 sin2
x 5. y

− 3y
+ 2y = cos
e−x  6. y

− 5y
+ 6y = 8 sin2
4x In each of Problems 7 through 16, find the general solution using the method of undetermined coefficients. 7. y

− y
− 2y = 2x2 + 5 8. y

− y
− 6y = 8e2x 9. y

− 2y
+ 10y = 20x2 + 2x − 8 10. y

− 4y
+ 5y = 21e2x 11. y

− 6y
+ 8y = 3ex 12. y

+ 6y
+ 9y = 9 cos
3x 13. y

− 3y
+ 2y = 10 sin
x

21. x2 y

− 5xy
+ 8y = 3x 4 22. x2 y

+ 3xy
+ y = x 23. x2 y

+ xy
+ 4y = sin
2 ln
x 24. x2 y

+ 2xy
− 6y = x2 − 2 25. y

− 4y
+ 4y = e3x − 1 26. y

− y
− 2y = x In each of Problems 27 through 42, solve the initial value problem. 27. y

− 4y = −7e2x + x y
0 = 1 y

0 = 3 28. y

+ 4y
= 8 + 34 cos
x y
0 = 3 y

0 = 2 29. y

+ 8y
+ 12y = e−x + 7 y
0 = 1 y

0 = 0 30. y

− 3y
= 2e2x sin
x y
0 = 1 y

0 = 2 31. y

− 2y
− 8y = 10e−x + 8e2x  y
0 = 1 y

0 = 4 32. y

− 6y
+ 9y = 4e3x  y
0 = 1 y

0 = 2 33. y

− 5y
+ 6y = cos
2x y
0 = 0 y

0 = 4 34. y

− y
+ y = 1 y
1 = 4 y

1 = −2

14. y

− 4y
= 8x2 + 2e3x

35. y

− 8y
+ 2y = e−x  y
−1 = 5 y

−1 = 2

15. y

− 4y
+ 13y = 3e2x − 5e3x

36. y

+ 6y
+ 9y = − cos
x y
0 = 1 y

0 = −6

16. y

− 2y
+ y = 3x + 25 sin
3x

37. y

− y = 5 sin2
x y
0 = 2 y

0 = −4

In each of Problems 17 through 26, find the general solution of the differential equation, using any method. 17. y

− y
− 2y = e2x 18. x2 y

+ 5xy
− 12y = ln
x 19. y

+ y
− 6y = x 20. y

− y
− 12y = 2 sinh2
x

2.7

93

38. y

+ y = tan
x y
0 = 4 y

0 = 3 39. x2 y

− 6y = 8x2  y
1 = 1 y

1 = 0 40. x2 y

+ 7xy
+ 9y = 27 ln
x y
1 = 1 y

1 = −4 41. x2 y

− 2xy
+ 2y = 10 sin
ln
x; y
1 = 3 y

1 = 0 42. x2 y

− 4xy
+ 6y = x4 ex  y
2 = 2 y

2 = 7

Application of Second-Order Differential Equations to a Mechanical System Envision a spring of natural (unstretched) length L and spring constant k. This constant quantifies the “stiffness” of the spring. The spring is suspended vertically. An object of mass m is attached at the lower end, stretching the spring d units past its rest length. The object comes to rest in its equilibrium position. It is then displaced vertically a distance y0 units (up or down), and released, possibly with an initial velocity (Figure 2.4). We want to construct a mathematical model allowing us to analyze the motion of the object. Let y
t be the displacement of the object from the equilibrium position at time t. As a convenience, take this equilibrium position to be y = 0. Choose down as the positive direction. Both of these choices are arbitrary.

94

CHAPTER 2

Second-Order Differential Equations (a) Unstretched

(b) Static equilibrium

(c) System in motion

L

d y0

m y m FIGURE 2.4

Mass/spring system.

Now consider the forces acting on the object. Gravity pulls it downward with a force of magnitude mg. By Hooke’s law, the force the spring exerts on the object has magnitude ky. At the equilibrium position, the force of the spring is −kd, negative because it acts upward. If the object is pulled downward a distance y from this position, an additional force −ky is exerted on it. Thus, the total force on the object due to the spring is −kd − ky The total force due to gravity and the spring is mg − kd − ky Since at the equilibrium point (y = 0) this force is zero, then mg = kd. The net force acting on the object due to gravity and the spring is therefore just −ky. Finally, there are forces tending to retard or damp out the motion. These include air resistance or viscosity of the medium if the object is suspended in some fluid such as oil. A standard assumption, arising from experiment, is that the retarding forces have magnitude proportional to the velocity y
. Thus, for some constant c called the damping constant, the retarding forces have magnitude cy
. The total force acting on the object due to gravity, damping and the spring itself therefore have magnitude −ky − cy
 Finally, there may be a driving force of magnitude f
t on the object. Now the total external force acting on the object has magnitude F = −ky − cy
+ f
t Assuming that the mass is constant, Newton’s second law of motion enables us to write my

= −ky − cy
+ f
t or y

+

c
k y + y = f
t m m

(2.14)

2.7 Application of Second-Order Differential Equations to a Mechanical System

95

This is the spring equation. We will analyze the motion described by solutions of this equation, under various conditions.

2.7.1

Unforced Motion

Suppose first that f
t = 0, so there is no driving force. Now the spring equation is y

+

c
k y + y=0 m m

2 +

k c  + = 0 m m

with characteristic equation

This has roots =−

c 1 √ 2 ± c − 4km 2m 2m

As we might expect, the general solution, hence the motion of the object, will depend on its mass, the amount of damping, and the stiffness of the spring. Consider the following cases. Case 1 c2 − 4km > 0 In this event, the characteristic equation has two real, distinct roots: 1 √ 2 1 √ 2 c c + − 1 = − c − 4km and 2 = − c − 4km 2m 2m 2m 2m The general solution of equation (2.14) in this case is y
t = c1 e1 t + c2 e2 t  Clearly 2 < 0. Since m and k are positive, c2 − 4km < c2 , so also. Therefore,

√

c2 − 4km < c and 1 is negative

lim y
t = 0

t→

regardless of initial conditions. In the case c2 − 4km > 0, the motion of the object decays to zero as time increases. This case is called overdamping, and it occurs when the square of the damping constant exceeds four times the product of the mass and spring constant.

EXAMPLE 2.27 Overdamping

Suppose c = 6, k = 5, and m = 1. Now the general solution is y
t = c1 e−t + c2 e−5t  Suppose, to be specific, the object was initially (at t = 0) drawn upward 4 units from the equilibrium position and released downward with a speed of 2 units per second. Then y
0 = −4 and y

0 = 2, and we obtain 1 y
t = e−t
−9 + e−4t  2 A graph of this solution is shown in Figure 2.5. What does the solution tell us about the motion? Since −9+e−4t < 0 for t > 0, then y
t < 0 and the object always remains above the equilibrium point. Its velocity y

t = e−t
9 − 5e−4t /2 decreases to zero as t increases, and y
t → 0 as t increases, so the object moves downward

96

CHAPTER 2

Second-Order Differential Equations y

y x 0

2

4

6

8

10

0

2

x 4

6

8

10


1


1


2


2


3


3


4


4

FIGURE 2.5 An example of overdamped motion, no driving force.

FIGURE 2.6 An example of critical damped motion, no driving force.

toward equilibrium with ever decreasing velocity, approaching closer to but never reaching the equilibrium point, and never coming to rest. Case 2 c2 − 4km = 0 Now the general solution of the spring equation (2.14) is y
t =
c1 + c2 te−ct/2m  This case is called critical damping. While y
t → 0 as t → , as in the overdamping case, we will see an important difference between critical and overdamping.

EXAMPLE 2.28

Let c = 2 and k = m = 1. Now y
t =
c1 + c2 te−t . Suppose the object is initially pulled up four units above the equilibrium position and then pushed downward with a speed of 5 units per second. Then y
0 = −4 and y

0 = 5, so y
t =
−4 + te−t  Observe that y
4 = 0, so, unlike the what we saw with overdamping, the object actually reaches the equilibrium position, four seconds after it was released, and then passes through it. In fact, y
t reaches its maximum when t = 5 seconds, and this maximum value is y
5 = e−5 , about 0007 unit below the equilibrium point. The velocity y

t =
5 − te−t is negative for t > 5, so the object’s velocity decreases after this 5-second point. Since y
t → 0 as t → , the object moves with decreasing velocity back toward the equilibrium point as time increases. Figure 2.6 shows a graph of the displacement function in this case. In general, when critical damping occurs, the object either passes through the equilibrium point exactly once, as just seen, or never reaches it at all, depending on the initial conditions. Case 3 c2 − 4km < 0 Now the spring constant and mass together are sufficiently large that c2 < 4km, and the damping is less dominant. This case is called underdamping. The general solution now is y
t = e−ct/2m c1 cos
t + c2 sin
t in which =

1 √ 4km − c2  2m

2.7 Application of Second-Order Differential Equations to a Mechanical System

97

Because c and m are positive, y
t → 0 as t → . However, now the motion is oscillatory because of the sine and cosine terms in the solution. The motion is not, however, periodic, because of the exponential factor, which causes the amplitude of the oscillations to decay to zero as time increases.

EXAMPLE 2.29

Suppose c = k = 2 and m = 1. Now the general solution is y
t = e−t c1 cos
t + c2 sin
t Suppose the object is driven downward from a point three units above equilibrium, with an initial speed of two units per second. Then y
0 = −3 and y

0 = 2 and the solution is y
t = −e−t
3 cos
t + sin
t The behavior of this solution is more easily visualized if we write it in phase angle form. We want to choose C and  so that 3 cos
t + sin
t = C cos
t +  For this, we need 3 cos
t + sin
t = C cos
t cos
 − C sin
t sin
 so C cos
 = 3

and

C sin
 = −1

Then C sin
 1 = tan
 = −  C cos
 3 so

    1 1  = tan−1 − = − tan−1  3 3

To solve for C, write C 2 cos2
 + C 2 sin2
 = C 2 = 32 + 12 = 10 √ so C = 10. Now we can write the solution as √ y
t = 10e−t cos
t − tan−1
1/3 The √ graph is therefore√a cosine curve with decaying amplitude, squashed between graphs of y = 10e−t and y = − 10e−t . The solution is shown in Figure 2.7, with these two exponential functions shown as reference curves. Because of the oscillatory cosine term, the object passes back and forth through the equilibrium point. In fact, it passes through equilibrium exactly when y
t = 0, or   1 2n + 1  t = tan−1 + 3 2 for n = 0 1 2 3 · · ·  In theory, the object oscillates through the equilibrium infinitely often in this underdamping case, although the amplitudes of the oscillations decrease to zero as time increases.

98

CHAPTER 2

Second-Order Differential Equations

3

10 e
t

2 1
1
2
3

0

4

2

8

6

t


10 e
t

FIGURE 2.7

An example of underdamped motion, no

driving force.

2.7.2

Forced Motion

Now suppose an external driving force of magnitude f
t acts on the object. Of course, different forces will cause different kinds of motion. As an illustration, we will analyze the motion under the influence of a periodic driving force f
t = A cos
t, with A and  positive constants. Now the spring equation is y+

A c
k y + y = cos
t m m m

(2.15)

We know how to solve this nonhomogeneous linear equation. Begin by finding a particular solution, using the method of undetermined coefficients. Attempt a solution yp
x = a cos
t + b sin
t Substitution of this into equation (2.15) and rearrangement of terms yields



k A k bc ac +a − −b −a2 + cos
t = b2 + sin
t m m m m m Since sin
t and cos
t are not constant multiples of each other, the only way this can be true for all t ≥ 0 is for the coefficient on each side of the equation to be zero. Therefore −a2 +

k A bc +a − = 0 m m m

and b2 +

k ac − b = 0 m m

Solve these for a and b, keeping in mind that A, c, k, and m are given. We get a= Let 0 = given by



A
k − m2 
k − m2 2 + 2 c2

and

b=

Ac 
k − m2 2 + 2 c2

k/m. Then a particular solution of equation (2.15), for this forcing function, is yp
x =

mA
20 − 2  cos
t m2
20 − 2 2 + 2 c2 +

assuming that c = 0 or  = 0 .

Ac m2
20 − 2 2 + 2 c2

sin
t

(2.16)

2.7 Application of Second-Order Differential Equations to a Mechanical System

99

We will now examine some specific cases to get some insight into the motion with this forcing function. Overdamped Forced Motion Suppose c =√6, k = 5, and √ m = 1, as we had previously in the overdamping case. Suppose also that A = 6 5 and  = 5. If the object is released from rest from the equilibrium position, then the displacement function satisfies the initial value problem √ √ y

+ 6y
+ 5y = 6 5 cos
5t y
0 = y

0 = 0 This problem has the unique solution √ y
t =

√ 5
−e−t + e−5t  + sin
5t 4

a graph of which is shown in Figure 2.8. As time increases, the exponential terms decrease to zero, exerting less influence on the motion,√while the sine term oscillates. Thus, as t increases, and the object moves up and down through the the solution tends to behave more like sin
5t √ equilibrium point, with approximate period 2/ 5. Contrast this with the overdamped motion with no forcing function, in which the object began above the equilibrium point and moved with decreasing velocity down toward it, but never reached it. Critically Damped Forced Motion Let c = 2 and m = k = 1. Suppose  = 1 and A = 2. Assume that the object is released from rest from the equilibrium position. Now the initial value problem for the position function is y

+ 2y
+ y = 2 cos
t

y
0 = y

0 = 0

with solution y
t = −te−t + sin
t A graph of this solution is shown in Figure 2.9. The exponential term exerts a significant influence at first, but decreases to zero as time increases. The term −te−t decreases to zero as √ 5 t increases, but not as quickly as the corresponding term 4
−e−t + e−5t  in the overdamping case. Nevertheless, after a while the motion settles into nearly (but not exactly, because −te−t is never actually zero for positive t) a sinusoidal motion back and forth through the equilibrium point. This is an example of critically damped forced motion.

1.0

1.0

0.5

0.5

0

5

10

15

20

t

0

5

10

15

20


0.5


0.5


1.0


1.0

FIGURE 2.8 of overdamped motion √ An example √ driven by 6 5 cos
5t.

FIGURE 2.9 An example of critical damped motion driven by 2 cos
t.

t

100

CHAPTER 2

Second-Order Differential Equations

1.0 0.5 0
0.5

5

10

15

20

t


1.0 FIGURE 2.10 example of underdamped motion √ An √ driven by 2 2 cos
2t.

√ √ Underdamped Forced Motion Suppose now that c = k = 2, m = 1,  = 2, and A = 2 2. Now c2 − 4km < 0, and we have underdamped motion, but this time with a forcing function. If the object is released from rest from the equilibrium position, then the initial value problem for the displacement function is √ √ y
0 = y

0 = 0 y

+ 2y
+ 2y = 2 2 cos
2t with solution

√ √ y
t = − 2e−t sin
t + sin
2t

Unlike the other two cases, the exponential factor in this solution has√a sin
t factor. Figure 2.10 shows a graph of this function. As time increases, the term − 2e−t sin
t becomes less influential and the motion settles √ nearly into an oscillation back and forth through the equilibrium point, with period nearly 2/ 2.

2.7.3

Resonance

In the absence of damping, an interesting phenomenon called resonance can occur. Suppose c = 0 but that there is still a periodic driving force f
t = A cos
t. Now the spring equation is A k y

+ y = cos
t m m From equation (2.16) with c = 0, this equation has general solution y
t = c1 cos
0 t + c2 sin
0 t +

A m
20 − 2 

cos
t

(2.17)

√ in which 0 = k/m. This number is called the natural frequency of the spring system, and is a function of the stiffness of the spring and mass of the object, while  is the input frequency and is contained in the driving force. This general solution assumes that the natural and input frequencies are different. Of course, the closer we choose the natural and input frequencies, the larger the amplitude of the cos
t term in the solution. Consider the case that the natural and input frequencies are the same. Now the differential equation is A k y = cos
0 t (2.18) m m and the function given by equation (2.17) is not a solution. To solve equation (2.18), first write the general solution yh of y

+
k/my = 0: y

+

yh
t = c1 cos
0 t + c2 sin
0 t

2.7 Application of Second-Order Differential Equations to a Mechanical System

101

For a particular solution of equation (2.18), we will proceed by the method of undetermined coefficients. Since the forcing function contains a term found in yh , we will attempt a particular solution of the form yp
t = at cos
0 t + bt sin
0 t Substitute this into equation (2.18) to obtain −2a0 sin
0 t + 2b0 cos
0 t =

A cos
0 t m

Thus choose A  m

a=0

and

yp
t =

A t sin
0 t 2m0

2b0 =

leading to the particular solution

The general solution of equation (2.18) is therefore y
t = c1 cos
0 t + c2 sin
0 t +

A t sin
0 t 2m0

This solution differs from that in the case  = 0 in the factor of t in yp
t. Because of this, solutions increase in amplitude as t increases. This phenomenon is called resonance. As a specific example, let c1 = c2 = 0 = 1 and A/2m = 1 to write the solution as y
t = cos
t + sin
t + t sin
t A graph of this function is shown in Figure 2.11, clearly revealing the increasing magnitude of the oscillations with time. While there is always some damping in the real world, if the damping constant is close to zero compared to other factors, such as the mass, and if the natural and input frequencies are (nearly) equal, then oscillations can build up to a sufficiently large amplitude to cause resonancelike behavior and damage a system. This can occur with soldiers marching in step across a bridge. If the cadence of the march (input frequency) is near enough to the natural frequency of the material of the bridge, vibrations can build up to dangerous levels. This occurred near Manchester, England, in 1831 when a column of soldiers marching across the Broughton Bridge caused it to collapse. More recently, the Tacoma Narrows Bridge in Washington experienced increasing oscillations driven by energy from the wind, causing it to whip about in sensational

20 10

0

5


10
20 FIGURE 2.11 Resonance.

10

15

20

t

102

CHAPTER 2

Second-Order Differential Equations

fashion before its collapse into the river. Videos of the wild thrashing about of the bridge are available in some libraries and engineering and science departments.

2.7.4

Beats

In the absence of damping, an oscillatory driving force can also cause a phenomenon called beats. Suppose  =  0 and consider y

+ 20 y =

A cos
0 t m

The Tacoma Narrows Bridge was completed in 1940 and stood as a new standard of combined artistry and functionality. The bridge soon became known for its tendency to sway in high winds, but no one suspected what was about to occur. On November 7, 1940, energy provided by unusually strong winds, coupled with a resonating effect in the bridge’s material and design, caused the oscillations in the bridge to be reinfored and build to dangerous levels. Soon, the twisting caused one side of the sidewalk to rise 28 feet above that of the other side. Concrete dropped out of the roadway, and a section of the suspension span completely rotated and fell away. Shortly thereafter, the entire center span collapsed into Puget Sound. This sensational construction failure motivated new mathematical treatments of vibration and wave phenomena in the design of bridges and other large structures. The forces that brought down this bridge are a more complicated version of the resonance phenomenon discussed in Section 2.7.3.

2.7 Application of Second-Order Differential Equations to a Mechanical System

103

1.0 0.5

0

5

10

15

20

25

30

t


0.5
1.0 FIGURE 2.12 Beats.

Assuming that the object is released from rest at the equilibrium position, then y
0 = y

0 = 0 and from equation (2.17) we have the solution y
t =

A

cos
t − cos
0 t m
20 − 2 

The behavior of this solution reveals itself more clearly if we write it as     1 1 2A

 sin + t sin − t  y
t = 2 0 2 0 m
20 − 2  This formulation reveals a periodic variation of amplitude in the solution, depending on the relative sizes of 0 +  and 0 − . It is this periodic variation of amplitude that is called a 1 beat. As a specific example, suppose 0 +  = 5 and 0 −  = , and the constants are chosen 2 so that 2A/ m
20 − 2  = 1. In this case, the displacement function is    5t t  y
t = sin sin 2 4 The beats are apparent in the graph of this solution in Figure 2.12.

2.7.5

Analogy with an Electrical Circuit

If a circuit contains a resistance R, inductance L, and capacitance C, and the electromotive force is E
t, then the impressed voltage is obtained as a sum of the voltage drops in the circuit: E
t = Li

t + Ri
t +

1 q
t C

Here i
t is the current at time t, and q
t is the charge. Since i = q
, we can write the second-order linear differential equation q

+

1 R
1 q + q = E L LC L

If R, L, and C are constant, this is a linear equation of the type we have solved for various choices of E
t. It is interesting to observe that this equation is of exactly the same form as the equation for the displacement of an object attached to a spring, which is y

+

1 c
k y + y = f
t m m m

CHAPTER 2

104

Second-Order Differential Equations

This means that solutions of one equation readily translate into solutions of the other and suggests the following equivalences between electrical and mechanical quantities: displacement function y
t ⇐⇒ charge q
t velocity y

t ⇐⇒ current i
t driving force f
t ⇐⇒ electromotive force E
t mass m ⇐⇒ inductance L damping constant c ⇐⇒ resistance R spring modulus k ⇐⇒ reciprocal 1/C of the capacitance

EXAMPLE 2.30

Consider the circuit of Figure 2.13, driven by a potential of E
t = 17 sin
2t volts. At time zero the current is zero and the charge on the capacitor is 1/2000 coulomb. The charge q
t on the capacitor for t > 0 is obtained by solving the initial value problem 10q

+ 120q
+ 1000q = 17 sin
2t q
0 =

1  q

0 = 0 2000

The solution is q
t =

1 −6t 1 e 7 cos
8t − sin
8t +

− cos
2t + 4 sin
2t 1500 240

0

1

3

2

4

5

t

6


0.01
0.02

120 

E(t)  17 sin(2t) volts

10
3 F
0.03

10 H FIGURE 2.13

FIGURE 2.14 Transient part of the current for the circuit of Figure 2.13.

0.03

0.03

0.02

0.02 0.01

0.01
0.01

0

2

4

6

8

10

t
0.01

0


0.02


0.02


0.03


0.03

FIGURE 2.15 Steady-state part of the current for the circuit of Figure 2.13.

Figure 2.13.

FIGURE 2.16

5

10

15

Current function for the circuit of

20

t

2.7 Application of Second-Order Differential Equations to a Mechanical System

105

The current can be calculated as 1 1

4 cos
2t + sin
2t i
t = q

t = − e−6t cos
8t + sin
8t + 30 120 The current is a sum of a transient part 1 −6t e cos
8t + sin
8t 30 named for the fact that it decays to zero as t increases, and a steady-state part −

1

4 cos
2t + sin
2t 120 The transient and steady-state parts are shown in Figures 2.14 and 2.15, and their sum, the current, is shown in Figure 2.16.

SECTION 2.7

PROBLEMS

1. The object of this problem is to gauge the relative effects of initial position and velocity on the motion in the unforced, overdamped case. Solve the initial value problems y

+ 4y
+ 2y = 0

y
0 = 5 y

0 = 0

y

+ 4y
+ 2y = 0

y
0 = 0 y

0 = 5

and

Graph the solutions on the same set of axes. What conclusions can be drawn from these solutions about the influence of initial position and velocity? 2. Repeat the experiment of Problem 1, except now use the critically damped unforced equation y

+ 4y
+ 4y = 0. 3. Repeat the experiment of Problem 1 for the underdamped unforced case y

+ 2y
+ 5y = 0.

8. y

+ 2y
+ 5y = 0 y
0 = A y

0 = 0; A has values 1 3 6 10 −4 and −7. 9. y

+ 2y
+ 5y = 0 y
0 = 0 y

0 = A; A has values 1 3 6 10 −4 and −7. 10. An object having mass 1 gram is attached to the lower end of a spring having spring modulus 29 dynes per centimeter. The bob is, in turn, adhered to a dashpot that imposes a damping force of 10v dynes, where v
t is the velocity at time t in centimeters per second. Determine the motion of the bob if it is pulled down 3 centimeters from equilibrium and then struck upward with a blow sufficient to impart a velocity of 1 centimeter per second. Graph the solution. Solve the problem when the initial velocity is, in turn, 2, 4, 7, and 12 centimeters per second. Graph these solutions on the same set of axes to visualize the influence of the initial velocity on the motion.

6. y

+ 4y
+ 4y = 0 y
0 = A y

0 = 0; A has values 1 3 6 10 −4 and −7.

11. An object having mass 1 kilogram is suspended from a spring having a spring constant of 24 newtons per meter. Attached to the object is a shock absorber, which induces a drag of 11v newtons (velocity is in meters per second). The system is set in motion by 25 lowering the bob centimeters and then striking it 3 hard enough to impart an upward velocity of 5 meters per second. Solve for and graph the displacement function. Obtain the solution for the cases that the bob is lowered, in turn, 12, 20, 30, and 45 centimeters, and graph the displacement functions for the five cases on the same set of axes to see the effect of the distance lowered.

7. y

+ 4y
+ 4y = 0 y
0 = 0 y

0 = A; A has values 1 3 6 10 −4 and −7.

12. When an 8-pound weight is suspended from a spring, it stretches the spring 2 inches. Determine the

Problems 4 through 9 explore the effects of changing the initial position or initial velocity on the motion of the bob. In each, use the same set of axes to graph the solution of the initial value problem for the given values of A and observe the effect that these changes cause in the solution. 4. y

+ 4y
+ 2y = 0 y
0 = A y

0 = 0; A has values 1 3 6 10 −4 and −7. 5. y

+ 4y
+ 2y = 0 y
0 = 0 y

0 = A; A has values 1 3 6 10 −4 and −7.

106

CHAPTER 2

Second-Order Differential Equations

equation of motion when an object with a mass of 7 kilograms is suspended from this spring, and the system is set in motion by striking the object an upward blow, imparting a velocity of 4 meters per second. 13. How many times can the bob pass through the equilibrium point in the case of overdamped motion? What condition can be placed on the initial displacement y
0 to guarantee that the bob never passes through equilibrium? 14. How many times can the bob pass through the equilibrium point in the case of critical damping? What condition can be placed on y
0 to ensure that the bob never passes through this position? How does the initial velocity influence whether the bob passes through the equilibrium position? 15. In underdamped motion, what effect does the damping constant c have on the frequency of the oscillations of motion? 16. Suppose y
0 = y

0 = 0. Determine the maximum displacement of the bob in the critically damped case, and show that the time at which this maximum occurs is independent of the initial displacement.

a fluid that imposes a drag of 2v pounds. The entire system is subjected to an external force 4 cos
t. Determine the value of  that maximizes the amplitude of the steady-state oscillation. What is this maximum amplitude? 21. Consider overdamped forced motion governed by y

+ 6y
+ 2y = 4 cos
3t. (a) Find the solution satisfying y
0 = 6 y

0 = 0. (b) Find the solution satisfying y
0 = 0 y

0 = 6. (c) Graph these solutions on the same set of axes to compare the effect of initial displacement with that of initial velocity. 22. Carry out the program of Problem 21 for the critically damped forced system governed by y

+ 4y
+ 4y = 4 cos
3t. 23. Carry out the program of Problem 21 for the underdamped forced system governed by y

+ y
+ 3y = 4 cos
3t. In each of Problems 24 through 27, use the information to find the current in the RLC circuit of Figure 2.17. Assume zero initial current and capacitor charge.

17. Suppose the acceleration of the bob on the spring at distance d from the equilibrium position is a. Prove √ that the period of the motion is 2 d/a in the case of undamped motion. 18. A mass m1 is attached to a spring and allowed to vibrate with undamped motion having period p. At some later time a second mass m2 is instantaneously fused with m1 . Prove that the new object, having mass m1 +m2 , exhibits simple harmonic motion with period p/ 1 + m2 /m1 . 19. Let y
t be the solution of y

+ 20 y =
A/m cos
t, with y
0 = y

0 = 0. Assuming that  = 0 , find lim→0 y
t. How does this limit compare with the solution of y

+ 20 y =
A/m cos
0 t, with y
0 = y

0 = 0? 20. A 16-pound weight is suspended from a spring, 8 stretching it feet. Then the weight is submerged in 11

R C L

FIGURE 2.17

RLC

circuit. 24. R = 200  L = 01 H C = 0006 F E
t = te−t volts 25. R = 400  L = 012 H C = 004 F E
t = 120 sin
20t volts 26. R = 150  L = 02 H C = 005 F, E
t = 1 − e−t volts 27. R = 450  L = 095 H C = 0007 F, E
t = e−t sin2
3t volts

CHAPTER

BASIC PROPERTIES INITIAL VALUE PROBLEMS USING THE LAPLACE TRANSFORM SHIFTING THEOREMS AND THE HEAVISIDE FUNCTION CONVOLUTION UNIT IMPULSES AND THE DIRAC DELTA FUNC-

3

The Laplace Transform

3.1

Definition and Basic Properties In mathematics, a transform is usually a device that converts one type of problem into another type, presumably easier to solve. The strategy is to solve the transformed problem, then transform back the other way to obtain the solution of the original problem. In the case of the Laplace transform, initial value problems are often converted to algebra problems, a process we can diagram as follows: initial value problem ⇓ algebra problem ⇓ solution of the algebra problem ⇓ solution of the initial value problem.

DEFINITION 3.1

Laplace Transform

The Laplace transform  f of f is a function defined by
e−st f
t dt  f
s = 0

for all s such that this integral converges. 107

108

CHAPTER 3

The Laplace Transform

The Laplace transform converts a function f to a new function called  f. Often we use t as the independent variable for f and s for the independent variable of  f. Thus, f
t is the function f evaluated at t, and  f
s is the function  f evaluated at s. It is often convenient to agree to use lowercase letters for a function put into the Laplace transform, and its upper case for the function that comes out. In this notation, F =  f

G =  g

H =  h

and so on.

EXAMPLE 3.1

Let f
t = eat , with a any real number. Then


 f
s = F
s =


0

k

e−st eat dt =






e
a−st dt

0



1
a−st = lim e e dt = lim k→ 0 k→ a − s

1
a−sk 1 = lim − e k→ a − s a−s =−

k


a−st

0

1 1 = a−s s−a

provided that a − s < 0, or s > a. The Laplace transform of f
t = eat is F
s = 1/
s − a, defined for s > a.

EXAMPLE 3.2

Let g
t = sin
t. Then  g
s = G
s = = lim




k→ 0

k






e−st sin
t dt

0

e−st sin
t dt

−ks

e cos k + se−ks sin k − 1 1 = lim −  = 2 k→ s2 + 1 s +1 G
s is defined for all s > 0. A Laplace transform is rarely computed by referring directly to the definition and integrating. Instead, we use tables of Laplace transforms of commonly used functions (such as Table 3.1) or computer software. We will also develop methods that are used to find the Laplace transform of a shifted or translated function, step functions, pulses, and various other functions that arise frequently in applications. The Laplace transform is linear, which means that constants factor through the transform, and the transform of a sum of functions is the sum of the transform of these functions.

3.1 Definition and Basic Properties

Table of Laplace Transforms of Functions

TA B L E 3 . 1 f
t 1.

1

2.

t

3.

tn
n = 1 2 3 · · · 

4.

1 √ t

5.

eat

6.

teat

7.

tn eat

8. 9.

1
eat − ebt  a−b 1
aeat − bebt  a−b

10.


c − beat +
a − cebt +
b − aect
a − b
b − c
c − a

11.

sin
at

12.

cos
at

13.

1 − cos
at

14.

at − sin
at

15.

sin
at − at cos
at

16

sin
at + at cos
at

17.

t sin
at

18.

t cos
at

19.

cos
at − cos
bt
b − a
b + a

20.

eat sin
bt

21.

eat cos
bt

22.

sinh
at

23.

cosh
at

24.

sin
atcosh
at − cos
atsinh
at

25.

sin
atsinh
at

F
s = f
t
s 1 s 1 s2 n! sn+1   s 1 s−a 1
s − a2 n!
s − an+1 1
s − a
s − b s
s − a
s − b 1
s − a
s − b
s − c a s 2 + a2 s s 2 + a2 a2 s
s2 + a2  a3 s2
s2 + a2  2a3 2
s + a2 2 2as2 2
s + a2 2 2as
s2 + a2 2
s2 − a2 
s2 + a2 2 s
s2 + a2 
s2 + b2  b
s − a2 + b2 s−a
s − a2 + b2 a s 2 − a2 s s 2 − a2 4a3 s4 + 4a4 2a2 s 4 s + 4a4

109

110

CHAPTER 3

The Laplace Transform

(continued)

TA B L E 3 . 1 f
t

F
s = f
t
s

26.

sinh
at − sin
at

2a3 s4 − a4

27.

cosh
at − cos
at

28.

1 √ eat
1 + 2at t

2a2 s s4 − a4 s
s − a3/2

29.

J0
at

√

30.

Jn
at

1 an

31.

√ J0
2 at 

32.

1 sin
at t

33.

2

1 − cos
at t

34.

2

1 − cosh
at t

1 s2 + a2 √

37.

  a 1 2 √ − aea t erfc √ t t   a 1 2 √ + aea t erf √ t t √ 2 ea t erf
a t 

38.

√ 2 ea t erfc
a t 

39.

erfc

35. 36.



40. 41. 42. 43. 44.

a √



2 t

1 2 √ e−a /4t t 1  
t + a √ 1 sin
2a t  t t f a t bt/a e f a

45.


t

46.


t − a

47.

Ln
t (Laguerre polynomial)

s2 + a2 − s √ s2 + a2

1 −a/s e s a tan−1 s  2  s + a2 ln 2 s  ln

s 2 − a2 s2



1 s+a √ s s − a2 a √ s
s − a2  1 √ √ s
s + a √

1 −a√s e s √ 1 √ e−a s s √ 1 √ eas erfc
as  s   a erf √ s aF
as aF
as − b e−s
1 − e−s  s e−as   1 s−1 n s s

n

3.1 Definition and Basic Properties

(continued)

TA B L E 3 . 1

48.

f
t

F
s = f
t
s

n! √ H2n
t
2n! t


1 − sn sn+1/2

(Hermite polynomial) 49.

√

−n! H2n+1
t 
2n + 1!


1 − sn sn+3/2

(Hermite polynomial) 50.

1 as2

triangular wave



1 − e−as 1 + e−as





=

 as 1 tanh 2 as 2

f (t) 1 a 51.

2a

3a

4a

t  as 1 tanh s 2

square wave

f (t) 1 a


1 52.

2a

3a

t

4a

e−as 1 − 2 as s
1 − e−as 

sawtooth wave

f (t) 1 a

2a

t

Operational Formulas f
t

F
s

af
t + bg
t

aF
s + bG
s




f
t f


sF
s − f
0+


t

sn F
s − sn−1 f
0 − · · · − f
n−1
0

f
 d

1 F
s s


n t

0

tf
t n

t f
t 1 f
t t

−F

s
−1n F
n
s
F
 d s

eat f
t

F
s − a

f
t − aH
t − a

e−as F
s

f
t +  = f
t (periodic)


 1 e−st f
t dt 1 − e−s 0

111

112

CHAPTER 3 THEOREM 3.1

The Laplace Transform Linearity of the Laplace Transform

Suppose  f
s and  g
s are defined for s > a, and  and  are real numbers. Then  f + g
s = F
s + G
s for s > a. Proof

 e−st f
t dt and 0 e−st g
t dt converge for s > a. Then
 f + g
s = e−st
f
t + g
t dt

By assumption,

 0

0

=




0

e−st f
t dt + 




0

e−st g
t dt = F
s + G
s

for s > a. This conclusion extends to any finite sum:  1 f1 + · · · + n fn 
s = 1 F1
s + · · · + n Fn
s for all s such that each Fj
s is defined.  Not every function has a Laplace transform, because 0 e−st f
t dt may not converge for any real values of s. We will consider conditions that can be placed on f to ensure that f has a Laplace transform.  k −st An obvious necessary  −st condition is that 0 e f
t dt must be defined for every k > 0, because  f
s = 0 e f
t dt. For this to occur, it is enough that f be piecewise continuous on 0 k for every positive number k. We will define this concept in general terms because it occurs in other contexts as well.

DEFINITION 3.2

Piecewise Continuity

f is piecewise continuous on a b if there are points a < t 1 < t2 < · · · < t n < b such that f is continuous on each open interval
a t1 ,
tj−1  tj , and
tn  b, and all of the following one-sided limits are finite: lim f
t lim f
t lim f
t and lim f
t

t→a+

t→tj −

t→tj +

t→b−

This means that f is continuous on a b except perhaps at finitely many points, at each of which f has finite one-sided limits from within the interval. The only discontinuities a piecewise continuous function f can experience on a b are finitely many jump discontinuities (gaps of finite width in the graph). Figure 3.1 shows typical jump discontinuities in a graph. For example, let ⎧ 2 t for 0 ≤ t ≤ 2 ⎪ ⎪ ⎪ ⎪ ⎪ ⎨ 2 at t = 2  f
t = ⎪ 1 for 2 < t ≤ 3 ⎪ ⎪ ⎪ ⎪ ⎩ −1 for 3 < t ≤ 4

3.1 Definition and Basic Properties

113

y(t) 4 3

y (t)

2 1 t1

t

t2

0
1

FIGURE 3.1 A function having jump discontinuities at t1 and t2 .

1

2

3

4

t

FIGURE 3.2

⎧ 2 t ⎪ ⎪ ⎪ ⎨ 2 f
t = ⎪ 1 ⎪ ⎪ ⎩ −1

if if if if

0≤t b. This line of reasoning suggests a set of conditions which are sufficient for a function to have a Laplace transform.

THEOREM 3.2

Existence of f

Suppose f is piecewise continuous on 0 k for every positive  k. Suppose also that there are numbers M and b, such that f
t ≤ Mebt for t ≥ 0. Then 0 e−st f
t dt converges for s > b, hence  f
s is defined for s > b. Many functions satisfy these conditions, including polynomials, sin
at, cos
at, eat , and others. The conditions of the theorem are sufficient, but not necessary for a function to have a Laplace transform. Consider, for example, f
t = t−1/2 for t > 0. This function is not piecewise

114

CHAPTER 3

The Laplace Transform

k continuous on any 0 k because limt→0+ t−1/2 = . Nevertheless, 0 e−st t−1/2 dt exists for every positive k and s > 0. Further,

2  f
s = e−st t−1/2 dt = 2 e−sx dx (let x = t1/2 ) 0

2
−z2 e dz =√ s 0   =  s

0

√ (let z = x s)

 2 √ in which we have used the fact (found in some standard integral tables) that 0 e−z dz = /2. Now revisit the flow chart at the start of this chapter. Taking the Laplace transform of a function is the first step in solving certain kinds of problems. The bottom of the flow chart suggests that at some point we must be able to go back the other way. After we find some function G
s, we will need to produce a function g whose Laplace transform is G. This is the process of taking an inverse Laplace transform.

DEFINITION 3.3

Given a function G, a function g such that  g = G is called an inverse Laplace transform of G. In this event, we write g = −1 G

For example,

−1

and



−1

1
t = eat s−a

1
t = sin
t s2 + 1

This inverse process is ambiguous because, given G, there will be be many functions whose Laplace transform is G. For example, we know that the Laplace transform of e−t is 1/
s + 1 for s > −1. However, if we change f
t at just one point, letting  e−t for t = 3 h
t = 0 for t = 3  −st  −st then 0 e f
t dt = 0 e h
t dt, and h has the same Laplace transform as f . In such a case, which one do we call the inverse Laplace transform of 1/
s + 1? One answer is provided by Lerch’s Theorem, which states that two continuous functions having the same Laplace transform must be equal. THEOREM 3.3

Lerch

Let f and g be continuous on 0  and suppose that  f =  g. Then f = g.

3.1 Definition and Basic Properties

115

In view of this, we will partially resolve the ambiguity in taking the inverse Laplace transform by agreeing that, given F
s, we seek a continuous f whose Laplace transform is F . If there is no continuous inverse transform function, then we simply have to make some agreement as to which of several possible candidates we will call −1 F. In applications, context will often make this choice obvious. Because of the linearity of the Laplace transform, its inverse is also linear. THEOREM 3.4

If −1 F = f and −1 G = g, and  and  are real numbers, then −1 F + G = f + g If Table 3.1 is used to find  f, look up f in the left column and read  f from the right column. For −1 F, look up F in the right column and match it with f in the left.

SECTION 3.1

PROBLEMS

In each of Problems 1 through 10, use the linearity of the Laplace transform, and Table 3.1, to find the Laplace transform of the function. 1. 2 sinh
t − 4 2. cos
t − sin
t 3. 4t sin
2t 4. t2 − 3t + 5 5. t − cos
5t 6. 2t2 e−3t − 4t + 1 7.
t + 42

5
s + 72 6 1 − 17. s − 4
s − 42

2 1 3 4 − 2+ 6 18. 4 s s s s 16.

Suppose that f
t is defined for all t ≥ 0. Then f is periodic with period T if f
t + T = f
t for all t ≥ 0. For example, sin
t has period 2. In Problems 19–22, assume that f has period T . 19. Show that

8. 3e−t + sin
6t

 f
s =

9. t − 3t + cos
4t 3

10. −3 cos
2t + 5 sin
4t In each of Problems 11 through 18, use the linearity of the inverse Laplace transform and Table 3.1 to find the (continuous) inverse Laplace transform of the function. 11. 12. 13. 14. 15.

−2 s + 16 4s s2 − 14 2s − 5 s2 + 16 3s + 17 s2 − 7 1 3 + s − 7 s2

20. Show that

n+1T nT





n+1T

e−st f
t dt

n=0 nT

e−st f
t dt = e−nsT




T

e−st f
t dt

0

21. From Problems 19 and 20, show that  
T  −nsT e e−st f
t dt  f
s = n=0

0



n 22. Use the geometric series n=0 r = 1/
1 − r for r < 1, together with the result of Problem 21, to show that
T 1  f
s = e−st f
t dt −sT 1−e 0

In each of Problems 23 through 30, a periodic function is given, sometimes by a graph. Find  f, using the result of Problem 22.

CHAPTER 3

116

The Laplace Transform

 5 23. f has period 6 and. f
t = 0

for 0 < t ≤ 3 for 3 < t ≤ 6

28. f has the graph of Figure 3.6.

f (t)

24. f
t = E sin
t, with E and  positive constants. (Here f has period /).

3

25. f has the graph of Figure 3.3.

2

f (t)

0

5

FIGURE 3.6

5 10

0

t

55 60

30 35

t

16 18

8 10

29. f has the graph of Figure 3.7.

FIGURE 3.3

f (t)

26. f has the graph of Figure 3.4.

h f (t) 2

a

0

3a

4a

5a

6a

7a

5a

6a

t

t

12

6

2a

FIGURE 3.7

FIGURE 3.4

30. f has the graph of Figure 3.8.

27. f has the graph of Figure 3.5.

f (t)

f (t) E sin(ωt)

h

E π

2π ω

ω

3π ω

t a

3a

4a

t

FIGURE 3.8

FIGURE 3.5

3.2

2a

Solution of Initial Value Problems Using the Laplace Transform The Laplace transform is a powerful tool for solving some kinds of initial value problems. The technique depends on the following fact about the Laplace transform of a derivative. THEOREM 3.5

Laplace Transform of a Derivative

Let f be continuous on 0  and suppose f
is piecewise continuous on 0 k for every positive k. Suppose also that limk→ e−sk f
k = 0 if s > 0. Then  f

s = sF
s − f
0

(3.1)

That is, the Laplace transform of the derivative of f is s times the Laplace transform of f at s, minus f at zero.

3.2 Solution of Initial Value Problems Using the Laplace Transform Proof

117

Begin with an integration by parts, with u = e−st and dv = f

t dt. For k > 0,
k
k e−st f

t dt = e−st f
tk0 − −se−st f
t dt 0

0

= e−sk f
k − f
0 + s




k

e−st f
t dt

0

Take the limit as k → and use the assumption that e−sk f
k → 0 to obtain


k
−sk −st  f 
s = lim e f
k − f
0 + s e f
t dt k→

= −f
0 + s


0

0



e−st f
t dt = −f
0 + sF
s

If f has a jump discontinuity at 0 (as occurs, for example, if f is an electromotive force that is switched on at time zero), then this conclusion can be amended to read  f

s = sF
s − f
0+ where f
0+ = lim f
t t→0+

is the right limit of f
t at 0. For problems involving differential equations of order 2 or higher, we need a higher derivative version of the theorem. Let f
j denote the jth derivative of f . As a notational convenience, we let f
0 = f . THEOREM 3.6

Laplace Transform of a Higher Derivative

Suppose f , f
, · · · , f n−1 are continuous on 0 , and f
n is piecewise continuous on 0 k for every positive k. Suppose also that limk→ e−sk f
j
k = 0 for s > 0 and for j = 1 2     n − 1. Then  f
n 
s = sn F
s − sn−1 f
0 − sn−2 f

0 − · · · − sf
n−2
0 − f
n−1
0

(3.2)

The second derivative case (n = 2) occurs sufficiently often that we will record it separately. Under the conditions of the theorem,  f


s = s2 F
s − sf
0 − f

0

(3.3)

We are now ready to use the Laplace transform to solve certain initial value problems.

EXAMPLE 3.3

Solve y
− 4y = 1 y
0 = 1 We know how to solve this problem, but we will use the Laplace transform to illustrate the technique. Write  y
s = Y
s. Take the Laplace transform of the differential equation, using the linearity of  and equation (3.1), with y
t in place of f
t:  y
− 4y
s =  y

s − 4 y
s 1 =
sY
s − y
0 − 4Y
s =  1
s =  s

118

CHAPTER 3

The Laplace Transform

Here we used the fact (from Table 3.1) that  1
s = 1/s for s > 0. Since y
0 = 1, we now have
s − 4Y
s = y
0 +

1 1 = 1+  s s

At this point we have an algebra problem to solve for Y
s, obtaining Y
s =

1 1 +
s − 4 s
s − 4

(note the flow chart at the beginning of this chapter). The solution of the initial value problem is

1 1 y = −1 Y = −1 + −1  s−4 s
s − 4 From entry 5 of Table 3.1, with a = 4, −1

1 = e4t  s−4

And from entry 8, with a = 0 and b = 4,

1 1 0t 1 −1  =
e − e4t  =
e4t − 1 s
s − 4 −4 4 The solution of the initial value problem is 1 y
t = e4t +
e4t − 1 4 5 4t 1 = e −  4 4 One feature of this Laplace transform technique is that the initial value given in the problem is naturally incorporated into the solution process through equation (3.1). We need not find the general solution first, then solve for the constant to satisfy the initial condition.

EXAMPLE 3.4

Solve y

+ 4y
+ 3y = et 

y
0 = 0 y

0 = 2

Apply  to the differential equation to get  y

 + 4 y
 + 3 y =  et  Now  y

 = s2 Y − sy
0 − y

0 = s2 Y − 2 and  y
 = sY − y
0 = sY Therefore, s2 Y − 2 + 4sY + 3Y =

1  s−1

3.2 Solution of Initial Value Problems Using the Laplace Transform

119

Solve for Y to obtain 2s − 1 
s − 1
s2 + 4s + 3

Y
s =

The solution is the inverse Laplace transform of this function. Some software will produce this inverse. If we want to use Table 3.1, we must use a partial fractions decomposition to write Y
s as a sum of simpler functions. Write Y
s = =

2s − 1
s − 1
s2 + 4s + 3 2s − 1 A B C = + + 
s − 1
s + 1
s + 3 s − 1 s + 1 s + 3

This equation can hold only if, for all s, A
s + 1
s + 3 + B
s − 1
s + 3 + C
s − 1
s + 1 = 2s − 1 Now choose values of s to simplify the task of determining A, B, and C. Let s = 1 to get 8A = 1, so A = 18 . Let s = −1 to get −4B = −3, so B = 43 . Choose s = −3 to get 8C = −7, so C = − 78 . Then Y
s =

1 1 3 1 7 1 + −  8 s−1 4 s+1 8 s+3

Now read from Table 3.1 that 1 3 7 y
t = et + e−t − e−3t  8 4 8 Again, the Laplace transform has converted an initial value problem to an algebra problem, incorporating the initial conditions into the algebraic manipulations. Once we obtain Y
s, the problem becomes one of inverting the transformed function to obtain y
t. Equation (3.1) has an interesting consequence that will be useful later. Under the conditions of the theorem, we know that  f
 = s f − f
0 Suppose f
t is defined by an integral, say f
t =




t

g
 d 0

Now f
0 = 0 and, assuming continuity of g, f

t = g
t. Then
t


 f  =  g = s g
 d  0

This means that 




t 0

1 g
 d =  g s

(3.4)

enabling us to take the Laplace transform of a function defined by an integral. We will use this equation later in dealing with circuits having discontinuous electromotive forces. Thus far we have illustrated a Laplace transform technique for solving initial value problems with constant coefficients. However, we could have solved the problems in these examples by other means. In the next three sections we will develop the machinery needed to apply the Laplace transform to problems that defy previous methods.

CHAPTER 3

120

The Laplace Transform

PROBLEMS

SECTION 3.2

11. Suppose f satisfies the hypotheses of Theorem 3.5, except for a jump discontinuity at 0. Show that  f

s = sF
s − f
0+, where f
0+ = limt→0+ f
t.

In each of Problems 1 through 10, use the Laplace transform to solve the initial value problem. 1. y
+ 4y = 1 y
0 = −3 2. y
− 9y = t y
0 = 5

12. Suppose f satisfies the hypotheses of Theorem 3.5, except for a jump discontinuity at a positive number c. Prove that

3. y
+ 4y = cos
t y
0 = 0 4. y
+ 2y = e−t  y
0 = 1

 f

s = sF
s − f
0 − e−cs f
c+ − f
c−

5. y
− 2y = 1 − t y
0 = 4 6. y

+ y = 1 y
0 = 6 y

0 = 0

where f
c− = limt→c− f
t.

7. y

− 4y
+ 4y = cos
t y
0 = 1 y

0 = −1

13. Suppose g is piecewise continuous on 0 k for every k > 0, and that there are numbers M, b, and a such that g
t ≤ Mebt for t ≥ a. Let  G = g. Show that


t 1
a 1 g
w dw
s = G
s − g
w dw  s s 0 0

8. y

+ 9y = t2  y
0 = y

0 = 0 9. y

+ 16y = 1 + t y
0 = −2 y

0 = 1 10. y

− 5y
+ 6y = e−t  y
0 = 0 y

0 = 2

3.3

Shifting Theorems and the Heaviside Function One point to developing the Laplace transform is to broaden the class of problems we are able to solve. Methods of Chapters 1 and 2 are primarily aimed at problems involving continuous functions. But many mathematical models deal with discontinuous processes (for example, switches thrown on and off in a circuit). For these, the Laplace transform is often effective, but we must learn more about representing discontinuous functions and applying both the transform and its inverse to them.

3.3.1

The First Shifting Theorem

We will show that the Laplace transform of eat f
t is nothing more than the Laplace transform of f
t, shifted a units to the right. This is achieved by replacing s by s − a in F
s to obtain F
s − a.

THEOREM 3.7

First Shifting Theorem, or Shifting in the s Variable

Let  f
s = F
s for s > b ≥ 0. Let a be any number. Then  eat f
t
s = F
s − a Proof

for s > a + b

Compute  eat f
t
s = =

for s − a > b, or s > a + b.






eat e−st f
s ds

0




0



e−
s−at f
t dt = F
s − a

3.3 Shifting Theorems and the Heaviside Function

121

EXAMPLE 3.5

We know from Table 3.1 that  cos
bt = s/
s2 + b2 . For the Laplace transform of eat cos
bt, replace s with s − a to get s−a 
s − a2 + b2

 eat cos
bt
s =

EXAMPLE 3.6

Since  t3  = 6/s4 , then  t3 e7t 
s =

6 
s − 74

The first shifting theorem suggests a corresponding formula for the inverse Laplace transform: If  f = F , then −1 F
s − a = eat f
t Sometimes it is convenient to write this result as −1 F
s − a = eat −1 F
s

(3.5)

EXAMPLE 3.7

Suppose we want to compute −1

4  s2 + 4s + 20

We will manipulate the quotient into a form to which we can apply the shifting theorem. Complete the square in the denominator to write 4 s2 + 4s + 20

=

4 
s + 22 + 16

Think of the quotient on the right as a function of s + 2: F
s + 2 =

4 
s + 22 + 16

This means we should choose F
s =

4 s2 + 16



Now the shifting theorem tells us that  e−2t sin
4t = F
s −
−2 = F
s + 2 =

4
s + 22 + 16

122

CHAPTER 3

The Laplace Transform

and therefore −1

4 = e−2t sin
4t
s + 22 + 16

EXAMPLE 3.8

Compute −1

3s − 1  s2 − 6s + 2

Again, begin with some manipulation into the form of a function of s − a for some a: 3s − 1 3s − 1 = s2 − 6s + 2
s − 32 − 7 =

8 3
s − 3 + = G
s − 3 + K
s − 3
s − 32 − 7
s − 32 − 7

if we choose G
s =

3s s2 − 7

and

K
s =

8 s2 − 7



Now apply equation (3.5) (in the second line) to write

3s − 1 −1  = −1 G
s − 3 + −1 K
s − 3 s2 − 6s + 2 = e3t −1 G
s + e3t −1 K
s

3s 8 + e3t −1 2 = e3t −1 2 s −7 s −7

 s  1 3t −1 3t −1 = 3e  + 8e  s2 − 7 s2 − 7 √ √ 8 = 3e3t cosh
7t + √ e3t sinh
7t 7

3.3.2

The Heaviside Function and Pulses

We will now lay the foundations for solving certain initial value problems having discontinuous forcing functions. To do this, we will use the Heaviside function. Recall that f has a jump discontinuity at a if limt→a− f
t and limt→a− f
t both exist and are finite, but unequal. Figure 3.9 shows a typical jump discontinuity. The magnitude of the jump discontinuity is the “width of the gap” in the graph at a. This width is     lim f
t − lim f
t   t→a−

t→a−

Functions with jump discontinuities can be treated very efficiently using the unit step function, or Heaviside function.

3.3 Shifting Theorems and the Heaviside Function

123

y

Magnitude of the jump discontinuity at a

x

a FIGURE 3.9

DEFINITION 3.4

Heaviside Function

The Heaviside function H is defined by H
t =



0 1

if t < 0 if t ≥ 0

Oliver Heaviside (1850–1925) was an English electrical engineer who did much to introduce Laplace transform methods into engineering practice. A graph of H is shown in Figure 3.10. It has a jump discontinuity of magnitude 1 at 0. The Heaviside function may be thought of as a flat switching function, “on” when t ≥ 0, where H
t = 1, and “off” when t < 0, where H
t = 0. We will use it to achieve a variety of effects, including switching functions on and off at different times, shifting functions along the axis, and combining functions with pulses. To begin this program, if a is any number, then H
t − a is the Heaviside function shifted a units to the right, as shown in Figure 3.11, since  0 if t < a H
t − a = 1 if t ≥ a H
t − a models a flat signal of magnitude 1, turned off until time t = a and then switched on. We can use H
t − a to achieve the effect of turning a given function g off until time t = a, at which time it is switched on. In particular,  0 if t < a H
t − ag
t = g
t if t ≥ a

H(t
a)

y

(0, 1)

1 x FIGURE 3.10

function H
t.

The Heaviside

a FIGURE 3.11 A shifted Heaviside

function.

t

124

CHAPTER 3

The Laplace Transform y

y

y  cos(t)

1.0

y  H(t
π)cos(t)

0.5

π
2

0

2

π 4

6

8

10

12

t

4

6

8

10

12

t


0.5
1.0 FIGURE 3.12

Comparison of y = cos
t and y = H
t −  cos
t.

is zero until time t = a, at which time it switches on g
t. To see this in a specific case, let g
t = cos
t for all t. Then  0 if t <  H
t − g
t = H
t −  cos
t = cos
t if t ≥  Graphs of cos
t and H
t − cos
t are shown in Figure 3.12 for comparison. We can also use the Heaviside function to describe a pulse.

DEFINITION 3.5

Pulse

A pulse is a function of the form k H
t − a − H
t − b in which a < b and k is a nonzero real number

This pulse function is graphed in Figure 3.13. It has value 0 if t < a (where H
t − a = H
t − b = 0), value 1 if a ≤ t < b (where H
t − a = 1 and H
t − b = 0), and value 0 if t ≥ b (where H
t − a = H
t − b = 1). Multiplying a function g by this pulse has the effect of leaving g
t switched off until time a. The function is then turned on until time b, when it is switched off again. For example, let g
t = et . Then ⎧ ⎪ ⎨0 if t < 1

H
t − 1 − H
t − 2et = et if 1 ≤ t < 2 ⎪ ⎩ 0 if t ≥ 2 Figure 3.14 shows a graph of this function. Next consider shifted functions of the form H
t −ag
t −a. If t < a, the g
t −aH
t −a = 0 because H
t − a = 0. If t ≥ a, then H
t − a = 1 and H
t − ag
t − a = g
t − a, which is g
t shifted a units to the right. Thus the graph of H
t − ag
t − a is zero along the horizontal axis until t = a, and for t ≥ a is the graph of g
t for t ≥ 0, shifted a units to the right to begin at a instead of 0.

3.3 Shifting Theorems and the Heaviside Function

125

f (t) 8 7 6 5 4 3 2 1

1 t a FIGURE 3.13

t 0

b

Pulse function H
t − a − H
t − b.

0.5

1.0

1.5

2.0

2.5

3.0

FIGURE 3.14 Graph of f
t = H
t − 1 − H
t − 2et .

EXAMPLE 3.9

Consider g
t = t2 and a = 2. Figure 3.15 compares the graph of g with the graph of H
t − 2 g
t − 2. The graph of g is a familiar parabola. The graph of H
t − 2g
t − 2 is zero until time 2, then has the shape of the graph of t2 for t ≥ 0, but shifted 2 units to the right to start at t = 2. It is important to understand the difference between g
t, H
t −ag
t, and H
t −ag
t −a. Figure 3.16 shows graphs of these three functions for g
t = t2 and a = 3.

3.3.3

The Second Shifting Theorem

Sometimes H
t − ag
t − a is referred to as a shifted function, although it is more than that because this graph is also zero for t < a. The second shifting theorem deals with the Laplace transform of such a function. y

y

t

t

2

FIGURE 3.15 Comparison of y = t 2 and

y =
t − 22 H
t − 2.

y

y

y

t

t 3

t 3

FIGURE 3.16 Comparison of y = t 2 , y = t 2 H
t − 3, and y =
t − 32 H
t − 3.

126

CHAPTER 3 THEOREM 3.8

The Laplace Transform Second Shifting Theorem, or Shifting In the t Variable

Let  f
s = F
s for s > b. Then  H
t − af
t − a
s = e−as F
s for s > b. That is, we obtain the Laplace transform of H
t − af
t − a by multiplying the Laplace transform of f
t by e−as . Proof

Proceeding from the definition,  H
t − af
t − a
s = =





0 a

e−st H
t − af
t − a dt e−st f
t − a dt

because H
t − a = 0 for t < a, and H
t − a = 1 for t ≥ a. Now let w = t − a in the last integral to obtain
 H
t − af
t − a
s = e−s
a+w f
w dw 0

= e−as




0

e−sw f
w dw = e−as F
s

EXAMPLE 3.10

Suppose we want the Laplace transform of H
t −a. Write this as H
t −af
t −a, with f
t = 1 for all t. Since F
s = 1/s (from Table 3.1 or by direct computation from the definition), then 1  H
t − a
s = e−as  1
s = e−as  s

EXAMPLE 3.11

Compute  g, where g
t = 0 for 0 ≤ t < 2 and g
t = t2 + 1 for t ≥ 2. Since g
t is zero until time t = 2, and is then t2 + 1, we may write g
t = H
t − 2
t2 + 1. To apply the second shifting theorem, we must write g
t as a function, or perhaps sum of functions, of the form f
t − 2H
t − 2. This necessitates writing t2 + 1 as a sum of functions of t − 2. One way to do this is to expand t2 + 1 in a Taylor series about 2. In this simple case we can achieve the same result by algebraic manipulation: t2 + 1 =
t − 2 + 22 + 1 =
t − 22 + 4
t − 2 + 5 Then g
t =
t2 + 1H
t − 2 =
t − 22 H
t − 2 + 4
t − 2H
t − 2 + 5H
t − 2

3.3 Shifting Theorems and the Heaviside Function

127

Now we can apply the second shifting theorem:  g = 
t − 22 H
t − 2 + 4
t − 2H
t − 2 + 5 H
t − 2 = e−2s  t2  + 4e−2s  t + 5e−2s  1

2 4 5  = e−2s 3 + 2 + s s s As usual, any formula for the Laplace transform of a class of functions can also be read as a formula for an inverse Laplace transform. The inverse version of the second shifting theorem is: −1 e−as F
s
t = H
t − af
t − a

(3.6)

This enables us to compute the inverse Laplace transform of a known transformed function multipled by an exponential e−as .

EXAMPLE 3.12

Compute 

−1

se−3s  s2 + 4

The presence of the exponential factor suggests the use of equation (3.6). Concentrate on finding

s −1   s2 + 4 This inverse can be read directly from Table 3.1, and is f
t = cos
2t. Therefore −3s se −1 2
t = H
t − 3 cos
2
t − 3 s +4 We are now prepared to solve certain initial value problems involving discontinuous forcing functions.

EXAMPLE 3.13

Solve the initial value problem y

+ 4y = f
t in which

y
0 = y

0 = 0 

f
t =

0 t

for t < 3 for t ≥ 3

Because of the discontinuity in f , methods developed in Chapter 2 do not apply. First recognize that f
t = H
t − 3t

128

CHAPTER 3

The Laplace Transform

Apply the Laplace transform to the differential equation to get  y

 +  y = s2 Y
s − sy
0 − y

0 + 4Y
s =
s2 + 4Y
s =  H
t − 3t in which we have inserted the initial conditions y
0 = y

0 = 0. In order to use the second shifting theorem to compute  H
t − 3t, write  H
t − 3t =  H
t − 3
t − 3 + 3 =  H
t − 3
t − 3 + 3 H
t − 3 = e−3s  t + 3e−3s  1 =

1 −3s 3 −3s e + e  s2 s

We now have
s2 + 4Y =

1 −3s 3 −3s e + e  s2 s

The transform of the solution is Y
s =

3s + 1 s2
s2 + 4

e−3s 

The solution is within reach. We must take the inverse Laplace transform of Y
s. To do this, first use a partial fractions decomposition to write 3s + 1 s2
s2 + 4

e−3s =

3 1 −3s 3 s 11 1 1 e − 2 e−3s + 2 e−3s − 2 e−3s  4s 4 s +4 4s 4 s +4

Each term is an exponential times a function whose Laplace transform we know, and we can apply equation (3.6) to write 3 3 y
t = H
t − 3 − H
t − 3 cos
2
t − 3 4 4 1 1 1 + H
t − 3
t − 3 − H
t − 3 sin
2
t − 3 4 4 2 Because of the H
t − 3 factor in each term, this solution is zero until time t = 3, and we may write  0 for t < 3 1 1 y
t = 3 3 − cos
2
t − 3 +
t − 3 − sin
2
t − 3 for t ≥ 3 4 4 4 8 or, upon combining terms,  0 for t < 3 y
t = 1

2t − 6 cos
2
t − 3 − sin
2
t − 3 for t ≥ 3 8 A graph of this solution is shown in Figure 3.17. In this example, it is interesting to observe that the solution is differentiable everywhere, even though the function f occurring in the differential equation had a jump discontinuity at 3. This behavior is typical of initial value problems having a discontinuous forcing function. If the differential equation has order n and  is a solution, then  and its first n − 1 derivatives will be continuous, while the nth derivative will have a jump discontinuity wherever f does, and these jump discontinuities will agree in magnitude with the corresponding jump discontinuities of f .

3.3 Shifting Theorems and the Heaviside Function

129

y 2.5 2.0 1.5 1.0 0.5 0

2

4

6

8

t

FIGURE 3.17 Solution of







y + 4y =

0 t

if 0 ≤ t < 3  if t ≤ 3


y
0 = y

0 = 0

Often we need to write a function having several jump discontinuities in terms of Heaviside functions in order to use the shifting theorems. Here is an example.

EXAMPLE 3.14

Let

⎧ ⎪ if t < 2 ⎨0 f
t = t − 1 if 2 ≤ t < 3 ⎪ ⎩ −4 if t ≥ 3

A graph of f is shown in Figure 3.18. There are jump discontinuities of magnitude 1 at t = 2 and magnitude 6 at t = 3. Think of f
t as consisting of two nonzero parts, the part that is t − 1 on 2 3 and the part that is −4 on 3 . We want to turn on t − 1 at time 2 and turn it off at time 3, then turn on −4 at time 3 and leave it on. The first effect is achieved by multiplying the pulse function H
t − 2 − H
t − 3 by t − 1. The second is achieved by multiplying H
t − 3 by 4. Therefore f
t = H
t − 2 − H
t − 3
t − 1 − 4H
t − 3 As a check, this gives f
t = 0 if t < 2 because all of the shifted Heaviside functions are zero for t < 2. For 2 ≤ t < 3, H
t − 2 = 1 but H
t − 3 = 0 so f
t = t − 1. And for t ≥ 3, H
t − 2 = H
t − 3 = 1, so f
t = −4.

3.3.4

Analysis of Electrical Circuits

The Heaviside function is important in many kinds of problems, including the analysis of electrical circuits, where we anticipate turning switches on and off. Here are two examples.

EXAMPLE 3.15

Suppose the capacitor in the circuit of Figure 3.19 initially has zero charge and that there is no initial current. At time t = 2 seconds, the switch is thrown from position B to A, held there for 1 second, then switched back to B. We want the output voltage Eout on the capacitor.

130

CHAPTER 3

The Laplace Transform

2 1 250,000 

t 2


1
2

3

4

5 A


3
4

B

10
6 F

E out

10 V

FIGURE ⎧ 3.18

⎪ ⎨

0 f
t = t − 1 ⎪ ⎩ −4

Graph of if t < 2 if 2 ≤ t < 3  if t ≥ 3

FIGURE 3.19

From the circuit diagram, the forcing function is zero until t = 2, then has value 10 volts until t = 3, and then is zero again. Thus E is the pulse function E
t = 10 H
t − 2 − H
t − 3 By Kirchhoff’s voltage law, Ri
t +

1 q
t = E
t C

or 250, 000q

t + 106 q
t = E
t We want to solve for q subject to the initial condition q
0 = 0. Apply the Laplace transform to the differential equation, incorporating the initial condition, to write 250000 sQ
t − q
0 + 106 Q
t = 250000sQ + 106 Q =  E
t Now  E
t
s = 10 H
t − 2
s − 10 H
t − 3
s 10 −2s 10 −3s e − e  s s We now have the following equation for Q: =

25
105 sQ
s + 106 Q
s =

10 −2s 10 −3s e − e s s

or Q
s = 4
10−5 

1 1 e−2s − 4
10−5 e−3s  s
s + 4 s
s + 4

Use a partial fractions decomposition to write



1 −2s 1 −3s −5 1 −2s −5 1 −3s − 10  Q
s = 10 e − e e − e s s+4 s s+4 By the second shifting theorem, −1



1 −2s e
t = H
t − 2 s

3.3 Shifting Theorems and the Heaviside Function

131

Eout (t)

E(t) 10

10 8 6 4

2 1

2

0

2

FIGURE 3.20

t

3

Input voltage for the circuit

of Figure 3.19.

and



−1

0

2

4

6

8

t

FIGURE 3.21 Output voltage for the circuit of Figure 3.19.

1 −2s = H
t − 2f
t − 2 e s+4

where f
t = −1 1/
s + 4 = e−4t . Thus

1 −2s −1 = H
t − 2e−4
t−2  e  s+4 The other two terms in Q
s are treated similarly, and we obtain q
t = 10−5 H
t − 2 − H
t − 2e−4
t−2  − 10−5 H
t − 3 − H
t − 3e−4
t−3  = 10−5 H
t − 2 1 − e−4
t−2  − 10−5 H
t − 3 1 − e−4
t−3  Finally, since the output voltage is Eout
t = 106 q
t, Eout
t = 10H
t − 2 1 − e−4
t−2  − 10H
t − 3 1 − e−4
t−3  The input and output voltages are graphed in Figures 3.20 and 3.21.

EXAMPLE 3.16

The circuit of Figure 3.22 has the roles of resistor and capacitor interchanged from the circuit of the preceding example. We want to know the output voltage i
tR at any time. The differential equation of the preceding example applies to this circuit, but now we are interested in the current. Since i = q
, then
25
105 i
t + 106 q
t = E
t

i
0 = q
0 = 0

10
6 F A 10 V

FIGURE 3.22

E out

B 2.5  10 5 

132

CHAPTER 3

The Laplace Transform Eout (t)

E(t)

10

10

5 0

2

3

4

5

6

t


5

2 1 0

1

t

2 3


10

FIGURE 3.23 Input voltage for the circuit of Figure 3.22.

FIGURE 3.24

Output voltage for the circuit

of Figure 3.22.

The strategy of eliminating q by differentiating and using i = q
does not apply here, because E
t is not differentiable. To eliminate q
t in the present case, write
t
t i
 d + q
0 = i
 d q
t = 0

0

We now have the following problem to solve for the current:
t
25
105 i
t + 106 i
 d = E
t 0

i
0 = 0

This is not a differential equation. Nevertheless, we have the means to solve it. Take the Laplace transform of the equation, using equation (3.4), to obtain 1
25
105 I
s + 106 I
s =  E
s s 1 1 = 10 e−2s − 10 e−3s  s s Here I =  i. Solve for I
s to get I
s = 4
10−5 

1 −2s 1 −3s e − 4
10−5  e  s+4 s+4

Take the inverse Laplace transform to obtain i
t = 4
10−5 H
t − 2e−4
t−2 − 4
10−5 H
t − 3e−4
t−3  The input and output voltages are graphed in Figures 3.23 and 3.24.

SECTION 3.3

PROBLEMS

In each of Problems 1 through 15, find the Laplace transform of the function. 1.
t3 − 3t + 2e−2t 2. e−3t
t − 2

 3. f
t =

1 cos
t

4. e4t t − cos
t

for 0 ≤ t < 7 for t ≥ 7

3.3 Shifting Theorems and the Heaviside Function  5. f
t =  6. f
t =

t 1 − 3t

for 0 ≤ t < 3 for t ≥ 3

2t − sin
t 0

for 0 ≤ t <  for t ≥ 

7. e−t 1 − t2 + sin
t  t2 for 0 ≤ t < 2 8. f
t = 1 − t − 3t2 for t ≥ 2  cos
t for 0 ≤ t < 2 9. f
t = 2 − sin
t for t ≥ 2 ⎧ ⎨−4 for 0 ≤ t < 1 0 for 1 ≤ t < 3 10. f
t = ⎩ −t e for t ≥ 3 11. te−2t cos
3t 12. et 1 − cosh
t  t − 2 for 0 ≤ t < 16 13. f
t = −1 for t ≥ 16  1 − cos
2t for 0 ≤ t < 3 14. f
t = 0 for t ≥ 3 15. e−5t
t4 + 2t2 + t In each of Problems 16 through 25, find the inverse Laplace transform of the function. 16.

17.

18.

1 s2 + 4s + 12 1 s2 − 4s + 5 1 −5s e s3

se−2s 19. 2 s +9 20.

3 −4s e s+2

21.

1 s2 + 6s + 7

s−4 22. 2 s − 8s + 10 23.

s+2 s2 + 6s + 1

24.

1 e−s
s − 53

25.

1 e−21s s
s2 + 16

133

t 26. Determine  e−2t 0 e2w cos
3w dw. Hint: Use the first shifting theorem. In each of Problems 27 through 32, solve the initial value problem by using the Laplace transform. 27.  y

+ 4y = f
t y
0 = 1 y

0 = 0, with f
t = 0 for 0 ≤ t < 4 3 for t ≥ 4 28. y

− 2y
− 3y = f
t y
0 = 1 y

0 = 0, with f
t = 0 for 0 ≤ t < 4 12 for t ≥ 4 29.  y
3 −8y = g
t y
0 = y

0 = y


0 = 0, with g
t = 0 for 0 ≤ t < 6 2 for t ≥ 6 30. y

+ 5y
+ 6y = f
t y
0 = y

0 = 0, with f
t = −2 for 0 ≤ t < 3 0 for t ≥ 3 31. y
3 − y

+ 4y
− 4y = f
t y
0 = y

0 = 0, 1 for 0 ≤ t < 5 y


0 = 1, with f
t = 2 for t ≥ 5 32. y

−4y
+4y = f
t y
0 = −2 y

0 = 1, with f
t = t for 0 ≤ t < 3 t + 2 for t ≥ 3 33. Calculate and graph the output voltage in the circuit of Figure 3.19, assuming that at time zero the capacitor is charged to a potential of 5 volts and the switch is opened at 0 and closed 5 seconds later. 34. Calculate and graph the output voltage in the RL circuit of Figure 3.25 if the current is initially zero and  0 for 0 ≤ t < 5 E
t = 2 for t ≥ 5

R E(t)

L

FIGURE 3.25

35. Solve for the current in the RL circuit of Problem  34 if the current is initially zero and E
t = k for 0 ≤ t < 5 0 for t ≥ 5

CHAPTER 3

134

The Laplace Transform

36. Solve for the current in the RL circuit of Problem  34 if the initial current is zero and E
t = 0 for 0 ≤ t < 4 Ae−t for t ≥ 4

f (t) h t a

37. Write the function graphed in Figure 3.26 in terms of the Heaviside function and find its Laplace transform.

b

c

FIGURE 3.28

f (t) 40. Solve for the current in the RL circuit of Figure 3.29 if the initial current is zero, E
t has period 4, and  10 for 0 ≤ t < 2 E
t = 0 for 2 ≤ t < 4

K t a

b

FIGURE 3.26

L E(t)

R

38. Write the function graphed in Figure 3.27 in terms of the Heaviside function and find its Laplace transform. FIGURE 3.29

f (t) Hint: See Problem 22 of Section 3.1 for the Laplace transform of a periodic function. You should find that I
s = F
s/
1 + e−2s  for some F
s. Use a geometric series to write  1 =
−1n e−2ns 1 + e−2s n=0

M t a

b

c

N FIGURE 3.27

to write I
s as an infinite series, then take the inverse transform term by term by using a shifting theorem. Graph the current for 0 ≤ t < 8.

39. Write the function graphed in Figure 3.28 in terms of the Heaviside function and find its Laplace transform.

3.4

Convolution In general the Laplace transform of the product of two functions is not the product of their transforms. There is, however, a special kind of product, denoted f ∗ g, called the convolution of f with g. Convolution has the feature that the transform of f ∗ g is the product of the transforms of f and g. This fact is called the convolution theorem.

DEFINITION 3.6

Convolution

If f and g are defined on 0 , then the convolution f ∗ g of f with g is the function defined by
t
f ∗ g
t = f
t − g
 d 0

for t ≥ 0.

3.4 Convolution THEOREM 3.9

135

Convolution Theorem

If f ∗ g is defined, then  f ∗ g =  f g Proof

Let F =  f and G =  g. Then F
sG
s = F
s






e−st g
t dt =

0






F
se−s g
 d

0

in which we changed the variable of integration to  and brought F
s within the integral. Now recall that e−s F
s =  H
t − f
t − 
s Substitute this into the integral for F
sG
s to get F
sG
s =




0

 H
t − f
t − 
sg
 d

(3.7)

But, from the definition of the Laplace transform,  H
t − f
t −  =




0

e−st H
t − f
t −  dt

Substitute this into equation (3.7) to get F
sG
s = =





0




0




0

0

e−st H
t − f
t −  dt g
 d

e−st g
H
t − f
t −  dt d

Now recall that H
t −  = 0 if 0 ≤ t < , while H
t −  = 1 if t ≥ . Therefore, F
sG
s =





0



e−st g
f
t −  dt d

Figure 3.30 shows the t plane. The last integration is over the shaded region, consisting of points
t  satisfying 0 ≤  ≤ t < . Reverse the order of integration to write F
sG
s = = =






t 0

0

e


0

−st




t 0

0

e−st g
f
t −  d

g
f
t −  d dt

e−st
f ∗ g
t dt =  f ∗ g
s

136

CHAPTER 3

The Laplace Transform

τ

tτ

t FIGURE 3.30

Therefore F
sG
s =  f ∗ g
s as we wanted to show. The inverse version of the convolution theorem is useful when we want to find the inverse transform of a function that is a product, and we know the inverse transform of each factor. THEOREM 3.10

Let −1 F = f and −1 G = g. Then −1 FG = f ∗ g

EXAMPLE 3.17

Compute 

−1

1  s
s − 42

We can do this several ways (a table, a program, a partial fractions decomposition). But we can also write



1 1 −1 1 −1 =  = −1 F
sG
s s
s − 42 s
s − 42 Now −1

1 = 1 = f
t s

and −1

1 = te4t = g
t
s − 42

Therefore, −1

1 = f
t ∗ g
t = 1 ∗ te4t s
s − 42
t 1 1 1 = e4 d = te4t − e4t +  4 16 16 0

The convolution operation is commutative.

3.4 Convolution

137

THEOREM 3.11

If f ∗ g is defined, so is g ∗ f , and f ∗ g = g ∗ f Proof

Let z = t −  in the integral defining the convolution to get
f ∗ g
t = =




t 0




0

t

f
t − g
 d f
zg
t − z
−1 dz =




t 0

f
zg
t − z dz =
g ∗ f
t

Commutativity can have practical importance, since the integral defining g ∗ f may be easier to evaluate than the integral defining f ∗ g in specific cases. Convolution can sometimes enable us to write solutions of problems that are stated in very general terms.

EXAMPLE 3.18

We will solve the problem y

− 2y
− 8y = f
t

y
0 = 1 y

0 = 0

Apply the Laplace transform, inserting the initial values, to obtain  y

− 2y
− 8y
s =
s2 Y
s − s − 2
sY
s − 1 − 8Y
s =  f
s = F
s Then
s2 − 2s − 8Y
s − s + 2 = F
s so Y
s =

s−2 1 F
s + 2  s2 − 2s − 8 s − 2s − 8

Use a partial fractions decomposition to write Y
s =

1 1 1 1 1 1 2 1 F
s − F
s + +  6 s−4 6 s+2 3 s−4 3 s+2

Then 1 1 1 2 y
t = e4t ∗ f
t − e−2t ∗ f
t + e4t + e−2t  6 6 3 3 This is the solution, for any function f having a convolution with e4t and e−2t . Convolution is also used to solve certain kinds of integral equations, in which the function to be determined occurs in an integral. We saw an example of this in solving for the current in Example 3.16.

138

CHAPTER 3

The Laplace Transform

EXAMPLE 3.19

Determine f such that f
t = 2t2 +




t 0

f
t − e− d

Recognize the integral on the right as the convolution of f with e−t . Thus the equation has the form f
t = 2t2 +
f ∗ e−t 
t Taking the Laplace transform of this equation yields F
s =

4 1  + F
s 3 s s+1

Then F
s =

4 4 + 4 3 s s

and from this we easily invert to obtain 2 f
t = 2t2 + t3  3

SECTION 3.4

PROBLEMS

In each of Problems 1 through 8, use the convolution theorem to compute the inverse Laplace transform of the function (even if another method would work). Wherever they occur, a and b are positive constants. 1.

1
s2 + 4
s2 − 4

1 e−2s s2 + 16 s 3. 2 2
s + a 
s2 + b2  2.

s2 4.
s − 3
s2 + 5

8.

2 s3
s2 + 5

In each of Problems 9 through 16, use the convolution theorem to write a formula for the solution of the initial value problem in terms of f
t. 9. y

− 5y
+ 6y = f
t y
0 = y

0 = 0 10. y

+ 10y
+ 24y = f
t y
0 = 1 y

0 = 0 11. y

− 8y
+ 12y = f
t y
0 = −3 y

0 = 2 12. y

− 4y
− 5y = f
t y
0 = 2 y

0 = 1 13. y

+ 9y = f
t y
0 = −1 y

0 = 1

5.

1 s
s2 + a2 2

6.

1 s4
s − 5

15. y
3 − y

− 4y
+ 4y = f
t y
0 = y

0 = 1, y


0 = 0

7.

1 e−4s s
s + 2

16. y
4 − 11y

+ 18y = f
t; y
0 = y

0 = y


0 = y
3
0 = 0

14. y

− k2 y = f
t y
0 = 2 y

0 = −4

3.5 Unit Impulses and the Dirac Delta Function

 t 23. f
t = e−3t et − 3 0 f
e3 d

In each of Problems 17 through 23, solve the integral equation. t 17. f
t = −1 + 0 f
t − e−3 d 18. f
t = −t +

t

19. f
t = e−t +

0

t 0

21. f
t = 3 +

t 0

25. Show by example that in general f ∗ 1 = f , where 1 denotes the function that is identically 1 for all t. Hint: Consider f
t = cos
t.

f
t −  d t 0

26. Use the convolution  t theorem to determine the Laplace transform of e−2t 0 e2w cos
3w dw.

f
t −  sin
 d

27. Use the convolution theorem to show that


t
w 1 −1  F
s
t = f
 d dw s2 0 0

f
 cos 2
t −  d

22. f
t = cos
t + e−2t

3.5

24. Use  t the convolution theorem to derive the formula  0 f
w dw
s =
1/sF
s. What assumptions are needed about f
t?

f
t −  sin
 d

20. f
t = −1 + t − 2

t 0

139

f
e2 d

Unit Impulses and the Dirac Delta Function Sometimes we encounter the concept of an impulse, which may be intuitively understood as a force of large magnitude applied over an instant of time. We can model an impulse as follows. For any positive number , consider the pulse  defined by 1 
t = H
t − H
t −   As shown in Figure 3.31, this is a pulse of magnitude 1/ and duration . By letting  approach zero, we obtain pulses of increasing magnitude over shorter time intervals. Dirac’s delta function is thought of as a pulse of “infinite magnitude” over an “infinitely short” duration, and is defined to be 
t = lim 
t →0+

This is not really a function in the conventional sense, but is a more general object called a distribution. Nevertheless, for historical reasons it continues to be referred to as the delta function. It is also named for the Nobel laureate physicist P.A.M. Dirac. The shifted delta function 
t − a is zero except for t = a, where it has its infinite spike. We can define the Laplace transform of the delta function as follows. Begin with 1 
t − a = H
t − a − H
t − a −  

1



t 0



FIGURE 3.31 Graph of


t − a.

140

CHAPTER 3

The Laplace Transform

Then

1 1 −as 1 −
a+s e−as
1 − e−s  = e − e   s s s

 
t − a = This suggests that we define

e−as
1 − e−s  = e−as  →0+ s

 
t − a = lim In particular, upon choosing a = 0 we have

 
t = 1 Thus we think of the delta function as having constant Laplace transform equal to 1. The following result is called the filtering property of the delta function. If at time a, a signal (function) is hit with an impulse, by multiplying it by 
t − a, and the resulting signal is summed over all positive time by integrating from zero to infinity, then we obtain exactly the signal value f
a. THEOREM 3.12

Filtering Property

Let a > 0 and let f be integrable on 0  and continuous at a. Then




f
t
t − a dt = f
a

0

Proof

First calculate


0

f
t
t − a dt =




0

1

H
t − a − H
t − a − f
t dt 

1
a+ = f
t dt  a

By the mean value theorem for integrals, there is some t between a and a +  such that


a+

a

f
t dt = f
t 

Then


0

f
t
t − a dt = f
t 

As  → 0+, a +  → a, so t → a and, by continuity, f
t  → f
a. Then
lim

→0+ 0



f
t
t − a dt = =

as we wanted to show.









0

f
t lim 
t − a dt →0+

0

f
t
t − a dt = lim f
t  = f
a →0+

3.5 Unit Impulses and the Dirac Delta Function

141

If we apply the filtering property to f
t = e−st , we get




0

e−st 
t − a dt = e−as 

consistent with the definition of the Laplace transform of the delta function. Further, if we change notation in the filtering property and write it as




0

f

 − t d = f
t

then we can recognize the convolution of f with  and read the last equation as f ∗  = f The delta function therefore acts as an identity for the “product” defined by the convolution of two functions. Here is an example of a boundary value problem involving the delta function.

EXAMPLE 3.20

Solve y

+ 2y
+ 2y = 
t − 3

y
0 = y

0 = 0

Apply the Laplace transform to the differential equation to get s2 Y
s + 2sY
s + 2Y
s = e−3s  hence Y
s =

e−3s s2 + 2s + 2



To find the inverse transform of the function on the right, first write Y
s =

1 e−3s 
s + 12 + 1

Now use both shifting theorems. Because −1 1/
s2 + 1 = sin
t, a shift in the s–variable gives us

1 −1 = e−t sin
t
s + 12 + 1 Now shift in the t–variable to obtain y
t = H
t − 3e−
t−3 sin
t − 3 A graph of this solution is shown in Figure 3.32. The solution is differentiable for t > 0, except that y

t has a jump discontinuity of magnitude 1 at t = 3. The magnitude of the jump is the coefficient of 
t − 3 in the differential equation.

142

CHAPTER 3

The Laplace Transform y(t) 0.30 0.25 0.20 0.15 0.10 0.05 0

2

FIGURE  3.32

y
t =

e

−
t−3

4

6

8

t

10

Graph of 0

if 0 ≤ t < 3

sin
t − 3

if t ≥ 3



The delta function may be used to study the behavior of a circuit that has been subjected to transients. These are generated during switching, and the high input voltages associated with them can create excessive current in the components, damaging the circuit. Transients can also be harmful because they contain a broad spectrum of frequencies. Introducing a transient into a circuit can therefore have the effect of forcing the circuit with a range of frequencies. If one of these is near the natural frequency of the system, resonance may occur, resulting in oscillations large enough to damage the system. For this reason, before a circuit is built, engineers sometimes use a delta function to model a transient and study its effect on the circuit.

EXAMPLE 3.21

Suppose, in the circuit of Figure 3.33, the current and charge on the capacitor are zero at time zero. We want to determine the output voltage response to a transient modeled by 
t The output voltage is q
t/C, so we will determine q
t. By Kirchhoff’s voltage law, Li
+ Ri +

1 q = i
+ 10i + 100q = 
t C

Since i = q
, q

+ 10q
+ 100q = 
t We assume initial conditions q
0 = q

0 = 0. Apply the Laplace transform to the differential equation and use the initial conditions to obtain s2 Q
s + 10sQ
s + 100Q
s = 1

E in (t)  δ(t)

FIGURE 3.33

1H

10 

0.01 F

3.5 Unit Impulses and the Dirac Delta Function

143

Then 1

Q
s =

s2 + 10s + 100



In order to invert this by using a shifting theorem, complete the square to write Q
s =

1 
s + 52 + 75

Since

√ 1 1 = √ sin
5 3t 2
s + 75 5 3



√ 1 1 = √ e−5t sin
5 3t 2
s + 5 + 75 5 3

−1 then −1

q
t =  The output voltage is

√ 1 20 q
t = 100q
t = √ e−5t sin
5 3t C 3 A graph of this output is shown in Figure 3.34. The circuit output displays damped oscillations at its natural frequency, even though it was not explicitly forced by oscillations of this frequency. If we wish, we can obtain the current by i
t = q

t. Eout (t) 5 4 3 2 1
1

0.5 0

1.0

1.5

2.0

2.5

3.0

t

FIGURE 3.34 Output of the circuit of

Figure 3.32.

SECTION 3.5

PROBLEMS

In each of Problems 1 through 5, solve the initial value problem and graph the solution. 1. y

+ 5y
+ 6y = 3
t − 2 − 4
t − 5; y
0 = y

0 = 0









2. y − 4y + 13y = 4
t − 3 y
0 = y
0 = 0

3. y
3 + 4y

+ 5y
+ 2y = 6
t; y
0 = y

0 = y


0 = 0 4. y

+ 16y
= 12
t − 5/8 y
0 = 3 y

0 = 0. 5. y

+ 5y
+ 6y = B
t y
0 = 3 y

0 = 0. Call the solution . What are 
0 and 

0? Using this

144

CHAPTER 3

The Laplace Transform

information, what physical phenomenon does the Dirac delta function model? 6. Suppose f is not continuous  at a, but limt→a+ f
t = f
a+ is finite. Prove that 0 f
t
t − a dt = f
a+.  7. Evaluate 0
sin
t/t
t − /6 dt. 2 8. Evaluate 0 t2 
t − 3 dt.  9. Evaluate 0 f
t
t − 2 dt, where ⎧ ⎪ ⎨t for 0 ≤ t < 2 f
t = t2 for t > 2 ⎪ ⎩ 5 for t = 2 10. It is sometimes convenient to consider 
t as the derivative of the Heaviside function H
t. Use the definitions of the derivative, the Heaviside function, and the delta function (as a limit of  ) to give a heuristic justification for this. 11. Use the idea that H

t = 
t from Problem 10 to determine the output voltage of the circuit of Example 3.16 by differentiating the relevant equation to obtain an equation in i rather than writing the charge as an integral. 12. If H

t = 
t, then  H

t
s = 1. Show that not all of the operational rules for the Laplace transform are compatible with this expression. Hint: Check to see whether H

t
s = s H
t
s − H
0+.

3.6

13. Evaluate 
t − a ∗ f
t. 14. An object of mass m is attached to the lower end of a spring of modulus k. Assume that there is no damping. Derive and solve an equation of motion for the position of the object at time t > 0, assuming that, at time zero, the object is pushed down from the equilibrium position with an initial velocity v0 . With what momentum does the object leave the equilibrium position? 15. Suppose an object of mass m is attached to the lower end of a spring having modulus k. Assume that there is no damping. Solve the equation of motion for the position of the object for any time t ≥ 0 if, at time zero, the weight is struck a downward blow of magnitude mv0 . How does the position of the object in Problem 14 compare with that of the object in this problem for any positive time? 16. A 2-pound weight is attached to the lower end of a spring, stretching it 83 inches. The weight is allowed to come to rest in the equilibrium position. At some later time, which is called time zero, the weight is struck a downward blow of magnitude 41 pound (an impulse). Assume that there is no damping in the system. Determine the velocity with which the weight leaves the equilibrium position as well as the frequency and magnitude of the resulting oscillations.

Laplace Transform Solution of Systems The Laplace transform can be of use in solving systems of equations involving derivatives and integrals.

EXAMPLE 3.22

Consider the system of differential equations and initial conditions for the functions x and y: x

− 2x
+ 3y
+ 2y = 4 2y
− x
+ 3y = 0 x
0 = x

0 = y
0 = 0 Begin by applying the Laplace transform to the differential equations, incorporating the initial conditions. We get 4 s2 X − 2sX + 3sY + 2Y =  s 2sY − sX + 3Y = 0

3.6 Laplace Transform Solution of Systems

145

Solve these equations for X
s and Y
s to get X
s =

2 4s + 6 and Y
s =  s2
s + 2
s − 1 s
s + 2
s − 1

A partial fractions decomposition yields X
s = −

1 1 1 10 1 71 −3 2 + + 2s s 6 s+2 3 s−1

and 2 1 1 1 1 +  Y
s = − + s 3 s+2 3 s−1 Upon applying the inverse Laplace transform, we obtain the solution 1 10 7 x
t = − − 3t + e−2t + et 2 6 3 and 1 2 y
t = −1 + e−2t + et  3 3 The analysis of mechanical and electrical systems having several components can lead to systems of differential equations that can be solved using the Laplace transform.

EXAMPLE 3.23

Consider the spring/mass system of Figure 3.35. Let x1 = x2 = 0 at the equilibrium position, where the weights are at rest. Choose the direction to the right as positive, and suppose the weights are at positions x1
t and x2
t at time t. By two applications of Hooke’s law, the restoring force on m1 is −k1 x1 + k2
x2 − x1  and that on m2 is −k2
x2 − x1  − k3 x2  By Newton’s second law of motion, m1 x1

= −
k1 + k2 x1 + k2 x2 + f1
t

k1

FIGURE 3.35

m1

k2

m2

k3

146

CHAPTER 3

The Laplace Transform

and m2 x2

= k2 x1 −
k2 + k3 x2 + f2
t These equations assume that damping is negligible, but allow for forcing functions acting on each mass. As a specific example, suppose m1 = m2 = 1 and k1 = k3 = 4 while k2 = 25 . Suppose f2
t = 0, so no external driving force acts on the second mass, while a force of magnitude f1
t = 2 1 − H
t − 3 acts on the first. This hits the first mass with a force of constant magnitude 2 for the first 3 seconds, then turns off. Now the system of equations for the displacement functions is x1

= −

13 5 x1 + x2 + 2 1 − H
t − 3 2 2 5 13 x2

= x1 − x2  2 2

If the masses are initially at rest at the equilibrium position, then x1
0 = x2
0 = x1

0 = x2

0 = 0 Apply the Laplace transform to each equation of the system to get 13 5 2
1 − e−3s  X1 + X2 +  2 2 s 5 13 s2 X2 = X1 − X2  2 2 s 2 X1 = −

Solve these to obtain X1
s =

  2 13 1 2
1 − e−3s  + s
s2 + 9
s2 + 4 2 s

and X2
s =

5 1
1 − e−3s 
s2 + 9
s2 + 4 s

In preparation for applying the inverse Laplace transform, use a partial fractions decomposition to write X1
s =

1 s 13 1 −3s 1 s 13 1 1 s 1 s − 2 − 2 − e + 2 e−3s + 2 e−3s 36 s 4 s + 4 9 s + 9 36 s 4 s +4 9 s +9

X2
s =

5 1 1 s 1 s 1 s 5 1 −3s 1 s − + − e + 2 e−3s − 2 e−3s 36 s 4 s2 + 4 9 s2 + 9 36 s 4 s +4 9 s +9

and

3.6 Laplace Transform Solution of Systems

147

Now it is routine to apply the inverse Laplace transform to obtain the solution x1
t =

1 13 1 − cos
2t − cos
3t 36 4 9

13 1 1 + − + cos
2
t − 3 − cos
3
t − 3 H
t − 3 36 4 9

x2
t =

1 5 − cos
2t + 19 cos
3t 36 4

5 1 1 + − + cos
2
t − 3 − cos
3
t − 3 H
t − 3 36 4 9

EXAMPLE 3.24

In the circuit of Figure 3.36, suppose the switch is closed at time zero. We want to know the current in each loop. Assume that both loop currents and the charges on the capacitors are initially zero. Apply Kirchhoff’s laws to each loop to get 40i1 + 120
q1 − q2  = 10 60i2 + 120q2 = 120
q1 − q2  Since i = q
, we can write q
t = 40i1 + 120 60i2 + 120




t 0




t 0

t 0

i
 d + q
0. Put into the two circuit equations, we get

i1
 − i2
 d + 120 q1
0 − q2
0 = 10

i2
d + 120q2
0 = 120




t 0

i1
 − i2
 d + 120 q1
0 − q2
0

Put q1
0 = q2
0 = 0 in this system to get 40i1 + 120 60i2 + 120





t 0 t 0

i1
 − i2
 d = 10 i2
 d = 120




t 0

i1
 − i2
 d

40 

60  1 120

10 V FIGURE 3.36

F

1 120

F

CHAPTER 3

148

The Laplace Transform

Apply the Laplace transform to each equation to get 120 120 10 I1 − I2 = s s s 120 120 120 I = I − I 60I2 + s 2 s 1 s 2 40I1 +

After some rearrangement, we have
s + 3I1 − 3I2 =

1 4

2I1 −
s + 4I2 = 0 Solve these to get I1
s =

3 1 1 1 s+4 = + 4
s + 1
s + 6 20 s + 1 10 s + 6

I2
s =

1 1 1 1 1 = −  2
s + 1
s + 6 10 s + 1 10 s + 6

and

Now use the inverse Laplace transform to find the solution i1
t =

SECTION 3.6

3 −t 1 −6t 1 1 e + e  i2
t = e−t − e−6t  20 10 10 10

PROBLEMS

In each of Problems 1 through 10, use the Laplace transform to solve the initial value problem for the system.

11. Use the Laplace transform to solve the system y1
− 2y2
+ 3y1 =0

1. x
− 2y
= 1 x
+ y − x = 0 x
0 = y
0 = 0

y1 − 4y2
+ 3y3
=t

2. 2x
− 3y + y
= 0 x
+ y
= t x
0 = y
0 = 0





y1 − 2y2
+ 3y3
= − 1 y1
0 = y2
0 = y3
0 = 0




3. x + 2y − y = 1 2x + y = 0 x
0 = y
0 = 0 4. x
+ y
− x = cos
2t x
+ 2y
= 0 x
0 = y
0 = 0

12. Solve for the currents in the circuit of Figure 3.37, assuming that the currents and charges are initially zero and that E
t = 2H
t − 4 − H
t − 5.

5. 3x
− y = 2t x
+ y
− y = 0 x
0 = y
0 = 0

2

6. x
+ 4y
− y = 0 x
+ 2y = e−t  x
0 = y
0 = 0 7. x
+ 2x − y
= 0 x
+ y + x = t2  x
0 = y
0 = 0 8. x
+ 4x − y = 0 x
+ y
= t x
0 = y
0 = 0

E(t)

10. x
+ 2y
− x = 0 4x
+ 3y
+ y = −6 x
0 = y
0 = 0

i1

FIGURE 3.37

i2 5H

3

9. x
+ y
+ x − y = 0 x
+ 2y
+ x = 1 x
0 = y
0 = 0

1

4

3.6 Laplace Transform Solution of Systems

149

highly varnished table. Show that, if stretched and released from rest, the masses oscillate with respect to each other with period

13. Solve for the currents in the circuit of Figure 3.37 if the currents and charges are initially zero and E
t = 1 − H
t − 4 sin
2
t − 4. 14. Solve for the displacement functions of the masses in the system of Figure 3.38. Neglect damping and assume zero initial displacements and velocities, and external forces f1
t = 2 and f2
t = 0. 15. Solve for the displacement functions in the system of Figure 3.38 if f1
t = 1 − H
t − 2 and f2
t = 0. Assume zero initial displacements and velocities.

 2

m1

m1 m2  k
m1 + m2 

m2

k

FIGURE 3.40

18. Solve for the currents in the circuit of Figure 3.41 if E
t = 5H
t − 2 and the initial currents are zero.

k1  6

m1  1

E(t)

k2  2

m2  1

16. Consider the system of Figure 3.39. Let M be subjected to a periodic driving force f
t = A sin
t. The masses are initially at rest in the equilibrium position. (a) Derive and solve the initial value problem for the displacement functions. (b) Show that, if m and k2 are chosen so that  =  k2 /m, then the mass m cancels the forced vibrations of M. In this case we call m a vibration absorber.

c1

y1

i2

10 

10 

20. Two tanks are connected by a series of pipes as shown in Figure 3.42. Tank 1 initially contains 60 gallons of brine in which 11 pounds of salt are dissolved. Tank 2 initially contains 7 pounds of salt dissolved in 18 gallons of brine. Beginning at time zero a mixture containing 16 pound of salt for each gallon of water is pumped into tank 1 at the rate of 2 gallons per minute, while salt water solutions are interchanged between the two tanks and also flow out of tank 2 at the rates shown in the diagram. Four minutes after time zero, salt is poured into tank 2 at the rate of 11 pounds per minute for a period of 2 minutes. Determine the amount of salt in each tank for any time t ≥ 0.

3 gal/ min

2 gal/ min k2

i1

19. Solve for the currents in the circuit of Figure 3.41 if E
t = 5
t − 1.

FIGURE 3.38

k1

30 H

FIGURE 3.41

k3  3

M

20 H

m 1 6

y2

lb/gal 1

2

FIGURE 3.39

17. Two objects of masses m1 and m2 are attached to opposite ends of a spring having spring constant k (Figure 3.40). The entire apparatus is placed on a

5 gal/ min FIGURE 3.42

2 gal/ min

150

CHAPTER 3

The Laplace Transform

21. Two tanks are connected by a series of pipes as shown in Figure 3.43. Tank 1 initially contains 200 gallons of brine in which 10 pounds of salt are dissolved. Tank 2 initially contains 5 pounds of salt dissolved in 100 gallons of water. Beginning at time zero, pure water is pumped into tank 1 at the rate of 3 gallons per minute, while brine solutions are interchanged between the tanks at the rates shown in the diagram. Three minutes after time zero, 5 pounds of salt are dumped into tank 2. Determine the amount of salt in each tank for any time t ≥ 0.

3.7

3 gal/ min

3 gal/ min

1

2 gal/ min

5 lb

2

4 gal/ min

1 gal/ min

FIGURE 3.43

Differential Equations with Polynomial Coefficients The Laplace transform can sometimes be used to solve linear differential equations having polynomials as coefficients. For this we need the fact that the Laplace transform of tf
t is the negative of the derivative of the Laplace transform of f
t.

THEOREM 3.13

Let  f
s = F
s for s > b and suppose that F is differentiable. Then  tf
t
s = −F

s for s > b. Proof

Differentiate under the integral sign to calculate
d d
−st
e−st f
t dt e f
t dt = ds 0 ds 0

−te−st f
t dt = e−st −tf
t dt =

F

s =

0

0

=  −tf
t
s and this is equivalent to the conclusion of the theorem. By applying this result n times, we reach the following.

COROLLARY 3.1

Let  f
s = F
s for s > b and let n be a positive integer. Suppose F is n times differentiable. Then, for s > b,  tn f
t
s =
−1n

dn F
s dsn

3.7 Differential Equations with Polynomial Coefficients

151

EXAMPLE 3.25

Consider the problem ty

+
4t − 2y
− 4y = 0

y
0 = 1

If we write this differential equation in the form y

+ p
ty
+ q
ty = 0, then we must choose p
t =
4t − 2/t, and this is not defined at t = 0, where the initial condition is given. This problem is not of the type for which we proved an existence/uniqueness theorem in Chapter 2. Further, we have only one initial condition. Nevertheless, we will look for functions satisfying the problem as stated. Apply the Laplace transform to the differential equation to get  ty

 + 4 ty

 − 2 y
 − 4 y = 0 Calculate the first three terms as follows. First,  ty

 = −

d d 2  y
 = −

s Y − sy
0 − y

0 ds ds

= −2sY − s2 Y
+ 1 because y
0 = 1 and y

0, though unknown, is constant and has zero derivative. Next, d


y  ds d = − sY − y
0 = −Y − sY
 ds

 ty
 = −

Finally,  y
 = sY − y
0 = sY − 1 The transform of the differential equation is therefore −2sY
− s2 Y + 1 − 4Y − 4sY
− 2sY + 2 − 4Y = 0 Then Y
+

3 4s + 8 Y=  s
s + 4 s
s + 4

This is a linear first-order differential equation, and we will find an integrating factor. First compute


4s + 8 ds = ln s2
s + 42  s
s + 4

Then eln
s

2
s+42 

= s2
s + 42

152

CHAPTER 3

The Laplace Transform

is an integrating factor. Multiply the differential equation by this factor to obtain s2
s + 42 Y
+
4s + 8s
s + 4Y = 3s
s + 4 or

s2
s + 42 Y
= 3s
s + 4 Integrate to get s2
s + 42 Y = s3 + 6s2 + C Then Y
s =

s 6 C + + 2  2 2
s + 4
s + 4 s
s + 42

Upon applying the inverse Laplace transform, we obtain y
t = e−4t + 2te−4t +

C

−1 + 2t + e−4t + 2te−4t  32

This function satisfies the differential equation and the condition y
0 = 1 for any real number C. This problem does not have a unique solution. When we applied the Laplace transform to a constant coefficient differential equation y

+ Ay
+ By = f
t, we obtained an algebraic expression for Y . In this example, with polynomials occurring as coefficients, we obtained a differential equation for Y because the process of computing the transform of tk y
t involves differentiating Y
s. In the next example, we will need the following fact.

THEOREM 3.14

Let f be piecewise continuous on 0 k for every positive number k, and suppose there are numbers M and b such that f
t ≤ Mebt for t ≥ 0. Let  f = F . Then lim F
s = 0

s→

Proof

Write 
 F
s = 

0


 e−st f
t dt ≤

M −
s−bt e = b−s

= 0



e−st Mebt dt

0

M →0 s−b

as s → . This result will enable us to solve the following initial value problem.

3.7 Differential Equations with Polynomial Coefficients

153

EXAMPLE 3.26

Suppose we want to solve y

+ 2ty
− 4y = 1

y
0 = y

0 = 0

Unlike the preceding example, this problem satisfies the hypotheses of the existence/uniqueness theorem in Chapter 2. Apply the Laplace transform to the differential equation to get 1 s2 Y
s − sy
0 − y

0 + 2 ty

s − 4Y
s =  s Now y
0 = y

0 = 0 and d

 y

s ds d

sY
s − y
0 = −Y
s − sY

s =− ds

 ty

s = −

We therefore have 1 s2 Y
s − 2Y
s − 2sY

s − 4Y
s =  s or  Y
+

 1 3 s − Y =− 2 s 2 2s

This is a linear first-order differential equation for Y . To find an integrating factor, first compute
3 s  1 − ds = 3 ln
s − s2  s 2 4 The exponential of this function, or s3 e−s

2 /4



is an integrating factor. Multiply the differential equation by this function to obtain
s3 e−s

2 /4

1 2 Y
= − se−s /4  2

Then s3 e−s

2 /4

Y = e−s

2 /4

+C

so Y
s =

1 C s2 /4 + e  s3 s3

CHAPTER 3

154

The Laplace Transform

We do not have any further initial conditions to determine C. However, in order to have lims→ Y
s = 0, we must choose C = 0. Then Y
s = 1/s3 so 1 y
t = t2  2

SECTION 3.7

PROBLEMS

Use the Laplace transform to solve each of Problems 1 through 10. 1. t2 y
− 2y = 2









6. y

+ 2ty
− 4y = 6 y
0 = 0 y

0 = 0 7. y

+ 8ty
= 0 y
0 = 4 y

0 = 0




2. y + 4ty − 4y = 0 y
0 = 0 y
0 = −7



5. ty

+
t − 1y
+ y = 0 y
0 = 0




3. y − 16ty + 32y = 14 y
0 = y
0 = 0 4. y

+ 8ty
− 8y = 0 y
0 = 0 y

0 = −4

8. y

− 4ty
+ 4y = 0 y
0 = 0 y

0 = 10 9. y

− 8ty
+ 16y = 3 y
0 = 0 y

0 = 0 10.
1 − ty

+ ty
− y = 0 y
0 = 3 y

0 = −1

CHAPTER

POWER SERIES SOLUTIONS OF INITIAL VALUE PROBLEMS SINGULAR POINTS AND THE METHOD OF FROBENIUSPOWER SERIES SOLUTIONS USING RECURRENCE RELATIONS SECOND SOLUTIONS AND LOGARITHM

4

Series Solutions

Sometimes we can find an explicit, closed form solution of a differential equation or initial value problem. This occurs with y
+ 2y = 1

y
0 = 3

which has the unique solution 1 y
x =
1 + 5e−2x  2 This solution is explicit, giving y
x as a function of x, and is in closed form because it is a finite algebraic combination of elementary functions (which are functions such as polynomials, trigonometric functions, and exponential functions). Sometimes standard methods do not yield a solution in closed form. For example, the problem y
+ ex y = x 2 

y
0 = 4

has the unique solution y
x = e−e

x




x 0

 2 ee d + 4e−e  

x

This solution is explicit, but it is not in closed form because of the integral. It is difficult to analyze this solution, or even to evaluate it at specific points. Sometimes a series solution is a good strategy for solving an initial value problem. Such a solution is explicit, giving y
x as an infinite series involving constants times powers of x. It may also reveal important information about the behavior of the solution—for example, whether it passes through the origin, whether it is an even or odd function, or whether the function is increasing or decreasing on a given interval. It may also be possible to make good approximations to function values from a series representation. We will begin with power series solutions for differential equations admitting such solutions. Following this, we will develop another kind of series for problems whose solutions do not have power series expansions about a particular point. This chapter assumes familiarity with basic facts about power series. 155

156

4.1

CHAPTER 4

Series Solutions

Power Series Solutions of Initial Value Problems Consider the linear first-order initial value problem y
+ p
xy = q
x

y
x0  = y0 

If p and q are continuous on an open interval I about x0 , we are guaranteed by Theorem 1.3 that this problem has a unique solution defined for all x in I. With a stronger condition on these coefficients, we can infer that the solution will have a stronger property, which we now define.

DEFINITION 4.1

Analytic Function

A function f is analytic at x0 if f
x has a power series representation in some open interval about x0 : f
x =



an
x − x0 n

n=0

in some interval
x0 − h x0 + h

For example, sin
x is analytic at 0, having the power series representation sin
x =


−1n 2n+1 x  n=0
2n + 1!

This series converges for all real x. Analyticity requires at least that f be infinitely differentiable at x0 , although this by itself is not sufficient for f to be analytic at x0 . We claim that, when the coefficients of an initial value problem are analytic, then the solution is as well. THEOREM 4.1

Let p and q be analytic at x0 . Then the initial value problem y
+ p
xy = q
x

y
x0  = y0

has a solution that is analytic at x0 . This means that an initial value problem whose coefficients are analytic at x0 has an analytic solution at x0 . This justifies attempting to expand the solution in a power series about x0 , where the initial condition is specified. This expansion has the form y
x =



an
x − x0 n 

n=0

in which an =

1
n y
x0  n!

(4.1)

4.1 Power Series Solutions of Initial Value Problems

157

One strategy to solve the initial value problem of the theorem is to use the differential equation and the initial condition to calculate these derivatives, hence obtain coefficients in the expansion (4.1) of the solution.

EXAMPLE 4.1

Consider again the problem y
+ ex y = x 2 

y
0 = 4

The theorem guarantees an analytic solution at 0: y
x =

 1
n y
0xn n=0 n!

= y
0 + y

0x +

1

1 y
0x2 + y
3
0x3 + · · · 2! 3!

We will know this series if we can determine the terms y
0, y

0, y


0, · · · . The initial condition gives us y
0 = 4. Put x = 0 into the differential equation to get y

0 + y
0 = 0 or y

0 + 4 = 0 Then y

0 = −4 Next determine y


0. Differentiate the differential equation to get y

+ ex y
+ ex y = 2x

(4.2)

and put x = 0 to get y


0 + y

0 + y
0 = 0 Then y


0 = −y

0 − y
0 = −
−4 − 4 = 0 Next we will find y
3
x. Differentiate equation (4.2) to get y
3 + 2ex y
+ ex y

+ ex y = 2 Then y
3
0 + 2y

0 + y


0 + y
0 = 2 or y
3
0 + 2
−4 + 4 = 2 Then y
3
0 = 6 Next differentiate equation (4.3): y
4 + 3ex y
+ 3ex y

+ ex y
3 + ex y = 0

(4.3)

158

CHAPTER 4

Series Solutions

Evaluate this at 0 to get y
4
0 + 3
−4 + 3
0 + 6 + 4 = 0 so y
4
0 = 2 At this point we have the first five terms of the Maclaurin expansion of the solution: 1 1 1 y
x = y
0 + y

0x + y


0x2 + y
3
0x3 + y
4
0x4 + · · · 2 6 24 1 1 = 4 − 4x + x3 + x4 + · · ·  6 12 By differentiating more times, we can write as many terms of this series as we want.

EXAMPLE 4.2

Consider the initial value problem y
+ sin
xy = 1 − x

y
 = −3

Since the initial condition is given at x = , we will seek terms in the Taylor expansion of the solution about . This series has the form 1 y
x =y
 + y


x −  + y



x − 2 2 +

1
3 1 y

x − 3 + y
4

x − 4 + · · ·  6 24

We know the first term, y
 = −3. From the differential equation, y

 = 1 −  + 3 sin
 = 1 −  Now differentiate the differential equation: y


x + cos
xy + sin
xy
= −1 Substitute x =  to get y


 −
−3 = −1 so y


 = −4 Next differentiate equation (4.4): y
3
x − sin
xy + 2 cos
xy
+ sin
xy

= 0 Substitute x =  to get y
3
 − 2
1 −  = 0 so y
3
 = 2
1 − 

(4.4)

4.1 Power Series Solutions of Initial Value Problems

159

Up to this point we have four terms of the expansion of the solution about : y
x = −3 +
1 − 
x −  −

4 2
1 − 
x − 2 +
x − 3 + · · · 2! 3!

1 = −3 +
1 − 
x −  − 2
x − 2 +
1 − 
x − 3 + · · ·  3 Again, with more work we can compute more terms. This method for generating a series solution of a first order linear initial value problem extends readily to second order problems, justified by the following theorem. THEOREM 4.2

Let p, q and f be analytic at x0 . Then the initial value problem y

+ p
xy
+ q
xy = f
x

y
x0  = A y

x0  = B

has a unique solution that is also analytic at x0 .

EXAMPLE 4.3

Solve y

− xy
+ ex y = 4

y
0 = 1 y

0 = 4

Methods from preceding chapters do not apply to this problem. Since −x, ex , and 4 are analytic at 0, the problem has a series solution expanded about 0. The solution has the form y
x = y
0 + y

0x +

1

1 y
0x2 + y
3
0x3 + · · ·  2! 3!

We already know the first two coefficients from the initial conditions. From the differential equation, y


0 = 4 − y
0 = 3 Now differentiate the differential equation to get y
3 − y
− xy

+ ex y + ex y
= 0 Then y
3
0 = y

0 − y
0 − y

0 = −1 Thus far we have four terms of the series solution about 0: 3 1 y
x = 1 + 4x + x2 − x3 + · · · 2 6 Although we have illustrated the series method for initial value problems, we can also use it to find general solutions.

CHAPTER 4

160

Series Solutions

EXAMPLE 4.4

We will find the general solution of y

+ cos
xy
+ 4y = 2x − 1 The idea is to think of this as an initial value problem, y

+ cos
xy
+ 4y = 2x − 1

y
0 = a y

0 = b

with a and b arbitrary (these will be the two arbitrary constants in the general solution). Now proceed as we have been doing. We will determine terms of a solution expanded about 0. The first two coefficients are a and b. For the coefficient of x2 , we find, from the differential equation y


0 = −y

0 − 4y
0 − 1 = −b − 4a − 1 Next, differentiate the differential equation: y
3 − sin
xy
+ cos
xy

+ 4y
= 2 so y
3
0 = −y


0 − 4y

0 + 2 = b + 4a + 1 − 4b + 2 = 4a − 3b + 3 Continuing in this way, we obtain (with details omitted) y
x = a + bx + +

−1 − 4a − b 2 3 + 4a − 3b 3 x + x 2 6

1 + 12a + 8b 4 −16 − 40a + b 5 x + x +···  24 120

In the next section, we will revisit power series solutions, but from a different perspective.

SECTION 4.1

PROBLEMS

In each of Problems 1 through 10, find the first five nonzero terms of the power series solution of the initial value problem, about the point where the initial conditions are given. 1. y

+ y
− xy = 0 y
0 = −2 y

0 = 0 2. y

+ 2xy
+
x − 1y = 0 y
0 = 1 y

0 = 2 3. y

− xy = 2x y
1 = 3 y

1 = 0 4. y

+ xy
= −1 + x y
2 = 1 y

2 = −4 5. y

−

1
1 y + y = 0 y
1 = 7 y

1 = 3 x2 x

8. y

+ y
− x4 y = sin
2x y
0 = 0 y

0 = −2 9. y

+

1
y − xy = 0 y
0 = y

0 = 1 x+2

1 10. y

− y
+ y = 1 y
4 = 0 y

4 = 2 x In each of Problems 11 through 20, find the first five nonzero terms of the Maclaurin expansion of the general solution. 11. y
+ sin
xy = −x 12. y
− x2 y = 1

6. y

+ x2 y = ex  y
0 = −2 y

0 = 7

13. y
+ xy = 1 − x + x2

7. y

− ex y
+ 2y = 1 y
0 = −3 y

0 = 1

14. y
− y = ln
x + 1

4.2 Power Series Solutions Using Recurrence Relations 15. y

+ xy = 0



In each of Problems 22 through 25, the initial value problem can be solved in closed form using methods from Chapters 1 and 2. Find this solution and expand it in a Maclaurin series. Then find the Maclaurin series solution using methods of Section 4.1. The two series should agree.




16. y − 2y + xy = 0 17. y

− x3 y = 1 18. y

+
1 − xy
+ 2xy = 0 19. y

+ y
− x2 y = 0

22. y

+ y = 1 y
0 = 0 y

0 = 0

20. y

− 8xy = 1 + 2x9

23. y
+ y = 2 y
0 = −1

21. Find the first five terms of the Maclaurin series solution of Airy’s equation y

+ xy = 0, satisfying y
0 = a y

0 = b.

4.2

161

24. y

+ 3y
+ 2y = x y
0 = 0 y

0 = 1 25. y

− 4y
+ 5y = 1 y
0 = −1 y

0 = 4

Power Series Solutions Using Recurrence Relations We have just seen one way to utilize the differential equation and initial conditions to generate terms of a series solution, expanded about the point where the initial conditions are specified. Another way to generate coefficients is to develop a recurrence relation, which allows us to produce coefficients once certain preceding ones are known. We will consider three examples of this method.

EXAMPLE 4.5

Consider y

+ x2 y = 0. Suppose we want a solution expanded about 0. Instead successive derivatives at 0, as we did before, now begin by substituting  of computing n y
x = n=0 an x into the differential equation. To do this, we need y
=



nan xn−1



and y

=

n=1

n
n − 1an xn−2 

n=2

Notice that the series for y
begins at n = 1, and that for y

at n = 2. Put these series into the differential equation to get y

+ x2 y =



n
n − 1an xn−2 +

n=2



an xn+2 = 0

(4.5)

n=0

Shift indices in both summations so that the power of x occurring in each series is the same. One way to do this is to write 

n
n − 1an xn−2 =

n=2




n + 2
n + 1an+2 xn

n=0

and  n=0

an xn+2 =



an−2 xn 

n=2

Using these series, we can write equation (4.5) as  n=0


n + 2
n + 1an+2 xn +

 n=2

an−2 xn = 0

162

CHAPTER 4

Series Solutions

We can combine the terms for n ≥ 2 under one summation and factor out the common xn (this was the reason for rewriting the series). When we do this, we must list the n = 0 and n = 1 terms of the first summation separately, or else we lose terms. We get 2
1a2 x0 + 3
2a3 x +




n + 2
n + 1an+2 + an−2 xn = 0

n=2

The only way for this series to be zero for all x in some open interval about 0 is for the coefficient of each power of x to be zero. Therefore, a2 = a3 = 0 and, for n = 2 3    ,
n + 2
n + 1an+2 + an−2 = 0 This implies that an+2 = −

1 a
n + 2
n + 1 n−2

for n = 2 3    

(4.6)

This is a recurrence relation for this differential equation. In this example, it gives an+2 in terms of an−2 for n = 2 3     Thus, we know a4 in terms of a0 , a5 in terms of a1 , a6 in terms of a2 , and so on. The form of the recurrence relation will vary with the differential equation, but it always gives coefficients in terms of one or more previously indexed ones. Using equation (4.6), we proceed: a4 = −

1 1 a = − a0 4
3 0 12

a5 = −

1 1 a = − a1 5
4 1 20

(by putting n = 2);

(by putting n = 3); a6 = −

1 a =0 6
5 2


because a2 = 0

a7 = −

1 a =0 7
6 3


because a3 = 0

a8 = −

1 1 a4 = a 8
7
56
12 0

a9 = −

1 1 a = a 9
8 5
72
20 1

and so on. The first few terms of the series solution expanded about 0 are 1 a x4 12 0 1 1 8 1 9 a1 x5 + 0x6 + 0x7 + x + x +··· − 20 672 1440     1 6 1 9 1 4 1 5 = a0 1 − x + x + · · · + a1 x − x + x +···  12 672 20 1440

y
x = a0 + a1 x + 0x2 + 0x3 −

This is actually the general solution, since a0 and a1 are arbitrary constants. Note that a0 = y
0 and a1 = y

0, so a solution is completely specified by giving y
0 and y

0.

4.2 Power Series Solutions Using Recurrence Relations

163

EXAMPLE 4.6

Consider the nonhomogeneous differential equation y

+ x2 y
+ 4y = 1 − x2  Attempt a solution y
x = 



n=0 an x

n

. Substitute this series into the differential equation to get

n
n − 1an xn−2 + x2

n=2

Then





nan xn−1 + 4

n=1

n
n − 1an xn−2 +

n=2





an x n = 1 − x 2 

n=0

nan xn+1 +

n=1



4an xn = 1 − x2 

(4.7)

n=0

Shift indices in the first and second summation so that the power of x occurring in each is xn : 

n
n − 1an xn−2 =

n=2

and




n + 2
n + 1an+2 xn

n=0



nan xn+1 =

n=1




n − 1an−1 xn 

n=2

Equation (4.7) becomes 


n + 2
n + 1an+2 xn +

n=0




n − 1an−1 xn +

n=2



4an xn = 1 − x2 

n=0

We can combine summations from n = 2 on, writing the n = 0 and n = 1 terms from the first and third summations separately. Then 2a2 x0 + 6a3 x + 4a0 x0 + 4a1 x +




n + 2
n + 1an+2 +
n − 1an−1 + 4an xn = 1 − x2 

n=2

For this to hold for all x in some interval about 0, the coefficient of xn on the left must match the coefficient of xn on the right. By matching these coefficients, we get: 2a2 + 4a0 = 1 (from x0 ), 6a3 + 4a1 = 0 (from x), 4
3a4 + a1 + 4a2 = −1 (from x2 ), and, for n ≥ 3,
n + 2
n + 1an+2 +
n − 1an−1 + 4an = 0

164

CHAPTER 4

Series Solutions

From these equations we get, in turn, 1 − 2a0  2 2 a3 = − a1  3 1 a4 =
−1 − a1 − 4a2  12   1 1 1 1 − 2a0 = − − a1 − 12 12 3 2

a2 =

1 1 2 = − + a0 − a1  4 3 12 and, for n = 3 4     an+2 = −

4an +
n − 1an−1 
n + 2
n + 1

This is the recurrence relation for this differential equation, and it enables us to determine an+2 if we know the two previous coefficients an and an−1 . With n = 3 we get   1 8 4a3 + 2a2 =− − a1 + 1 − 4a0 a5 = − 20 20 3 =−

2 1 1 + a0 + a1  20 5 15

With n = 4 the recurrence relation gives us a6 = − =

  1 1 1 8
4a4 + 3a3  = − −1 + a0 − a1 − 2a1 30 30 3 3

4 7 1 − a + a 30 45 0 90 1

Thus far we have six terms of the solution:   1 2 y
x = a0 + a1 x + − 2a0 x2 − a1 x3 2 3     1 1 2 1 1 2 4 + − + a0 − a1 x + − + a0 + a1 x5 4 3 12 20 5 15   4 1 7 − a + a x6 + · · ·  + 30 45 0 90 1

Using the recurrence relation, we can produce as many terms of this series as we wish. A recurrence relation is particularly suited to computer generation of coefficients. Because this recurrence relation specifies each an (for n ≥ 3) in terms of two preceding coefficients, it is called a two-term recurrence relation. It will give each an for n ≥ 3 in terms of a0 and a1 , which are arbitrary constants. Indeed, y
0 = a0 and y

0 = a1 , so assigning values to these constants uniquely determines the solution. Sometimes we must represent one or more coefficients as power series to apply the current method. This does not alter the basic idea of collecting coefficients of like powers of x and solving for the coefficients.

4.2 Power Series Solutions Using Recurrence Relations

165

EXAMPLE 4.7

Solve y

+ xy
− y = e3x  Each coefficient at 0, so we will series solution expanded about 0. is analytic  look nfor a power n 3x n Substitute y = a x and also e =
3 /n!x into the differential equation to get: n=0 n n=0 

n
n − 1an xn−2 +

n=2



nan xn −

n=1



an xn =

n=0

 3n n=0

n!

xn 

Shift indices in the first summation to write this equation as 


n + 2
n + 1an+2 xn +

n=0



nan xn −

n=1



an x n =

n=0

 3n n=0

n!

xn 

We can collect terms from n = 1 on under one summation, obtaining 


n + 2
n + 1an+2 +
n − 1an xn + 2a2 − a0 = 1 +

n=1

 3n n=1

n!

xn 

Equate coefficients of like powers of x on both sides of the equation to obtain 2a2 − a0 = 1 and, for n = 1 2    ,
n + 2
n + 1an+2 +
n − 1an =

3n  n!

This gives 1 a2 =
1 + a0  2 and, for n = 1 2     we have the one-term recurrence relation (in terms of one preceding coefficient) an+2 =


3n /n! +
1 − nan 
n + 2
n + 1

Using this relationship we can generate as many coefficients as we want in the solution series, in terms of the arbitrary constants a0 and a1 . The first few terms are 1 + a0 2 1 3 y
x a0 + a1 x + x + x 2 2     1 a0 7 1 57 a0 − + + x4 + x5 + x6 + · · ·  3 24 40 30 24 8

SECTION 4.2

PROBLEMS

In each of Problems 1 through 12, find the recurrence relation and use it to generate the first five terms of the Maclaurin series of the general solution.

1. y
− xy = 1 − x 2. y
− x3 y = 4

166

CHAPTER 4

Series Solutions

3. y
+
1 − x2 y = x

8. y

+ x2 y
+ 2y = 0

4. y

+ 2y
+ xy = 0

9. y

+
1 − xy
+ 2y = 1 − x2

5. y

− xy
+ y = 3

10. y

+ y
−
1 − x + x2 y = −5

6. y

+ xy
+ xy = 0

11. y
+ xy = cos
x

7. y

− x2 y
+ 2y = x

12. y

+ xy
= 1 − ex

4.3

Singular Points and the Method of Frobenius In this section we will consider the second-order linear differential equation P
xy

+ Q
xy
+ R
xy = F
x

(4.8)

If we can divide this equation by P
x and obtain an equation of the form y

+ p
xy
+ q
xy = f
x

(4.9)

with analytic coefficients in some open interval about x0 , then we can proceed to a power series solution of equation (4.9) by methods already developed, and thereby solve equation (4.8). In this case we call x0 an ordinary point of the differential equation. If, however, x0 is not an ordinary point, then this strategy fails and we must develop some new machinery.

DEFINITION 4.2

Ordinary and Singular Points

x0 is an ordinary point of equation (4.8) if P
x0  = 0 and Q
x/P
x, R
x/P
x, and F
x/P
x are analytic at x0 . x0 is a singular point of equation (4.8) if x0 is not an ordinary point.

Thus, x0 is a singular point if P
x0  = 0, or if any one of Q
x/P
x, R
x/P
x, or F
x/P
x fails to be analytic at x0 .

EXAMPLE 4.8

The differential equation x3
x − 22 y

+ 5
x + 2
x − 2y
+ 3x2 y = 0 has singular points at 0 and 2, because P
x = x3
x − 22 and P
0 = P
2 = 0. Every other real number is a regular point of this equation. In an interval about a singular point, solutions can exhibit behavior that is quite different from what we have seen in an interval about an ordinary point. In particular, the general solution of equation (4.8) may contain a logarithm term, which will tend toward in magnitude as x approaches x0 .

4.3 Singular Points and the Method of Frobenius

167

In order to seek some understanding of the behavior of solutions near a singular point, we will concentrate on the homogeneous equation P
xy

+ Q
xy
+ R
xy = 0

(4.10)

Once this case is understood, it does not add substantial further difficulty to consider the nonhomogeneous equation (4.8). Experience and research have shown that some singular points are “worse” than others, in the sense that the subtleties they bring to attempts at solution are deepened. We therefore distinguish two kinds of singular points.

DEFINITION 4.3

Regular and Irregular Singular Points

x0 is a regular singular point of equation (4.10) if x0 is a singular point, and the functions
x − x0 

Q
x P
x

and
x − x0 2

R
x P
x

are analytic at x0 . A singular point that is not regular is said to be an irregular singular point.

EXAMPLE 4.9

We have already noted that x3
x − 22 y

+ 5
x + 2
x − 2y
+ 3x2 y = 0 has singular points at 0 and 2. We will classify these singular points. In this example, P
x = x3
x − 22 , Q
x = 5
x + 2
x − 2 and R
x = 3x2 . First consider x0 = 0. Now
x − x0 

Q
x 5x
x + 2
x − 2 5 x+2 = = 2 P
x x3
x − 22 x x−2

is not defined at 0, hence is not analytic there. This is enough to conclude that 0 is an irregular singular point of this differential equation. Next let x0 = 2 and consider
x − 2

x+2 Q
x =5 3 P
x x

and
x − 22

R
x 3 =  P
x x

Both of these functions are analytic at 2. Therefore, 2 is a regular singular point of the differential equation. Suppose now that equation (4.10) has a regular singular point at x0 . Then there may be no solution as a power series about x0 . In this case we attempt to choose numbers cn and a number r so that  y
x = cn
x − x0 n+r (4.11) n=0

168

CHAPTER 4

Series Solutions

is a solution. This series is called a Frobenius series, and the strategy of attempting a solution of this form is called the method of Frobenius. A Frobenius series need not be a power series, since r may be negative or may be a noninteger. A Frobenius series “begins” with c0 xr , which is constant only if r = 0. Thus, in computing the derivative of the Frobenius series (4.11), we get y

x =




n + rcn
x − x0 n+r−1 

n=0

and this summation begins at zero again because the derivative of the n = 0 term need not be zero. Similarly, y


x =




n + r
n + r − 1cn
x − x0 n+r−2 

n=0

We will now illustrate the method of Frobenius.

EXAMPLE 4.10

We want to solve

 x2 y

+ x

   1 1 + 2x y
+ x − y = 0 2 2

It to show that 0 is a regular singular point. Substitute a Frobenius series y
x =  is routine n+r into the differential equation to get n=0 cn x 


n + r
n + r − 1cn xn+r +

n=0

 1 n=0

+

 n=0

cn xn+r+1 −

 1 n=0 2

2


n + rcn xn+r +



2
n + rcn xn+r+1

n=0

cn xn+r = 0

Shift indices in the third and fourth summations to write this equation as

 1 1 1 r r
r − 1c0 + c0 r − c0 x +
n + r
n + r − 1cn +
n + rcn 2 2 2 n=1

1 + 2
n + r − 1cn−1 + cn−1 − cn xn+r = 0 2 This equation will hold if the coefficient of each xn+r is zero. This gives us the equations 1 1 r
r − 1c0 + c0 r − c0 = 0 2 2

(4.12)

1 1
n + r
n + r − 1cn +
n + rcn + 2
n + r − 1cn−1 + cn−1 − cn = 0 2 2

(4.13)

and

for n = 1 2     Assuming that c0 = 0, an essential requirement in the method, equation (4.12) implies that 1 1 r
r − 1 + r − = 0 2 2

(4.14)

4.3 Singular Points and the Method of Frobenius

169

This is the indicial equation for this differential equation, and it determines the values of r we 1 can use. Solve it to obtain r1 = 1 and r2 = − . Equation (4.13) enables us to solve for cn in 2 terms of cn−1 to get the recurrence relation cn = −

1 + 2
n + r − 1 c 1 1 n−1
n + r
n + r − 1 +
n + r − 2 2

for n = 1 2     First put r = r1 = 1 into the recurrence relation to obtain 1 + 2n c cn = −  3 n−1 n n+ 2 Some of these coefficients are

for n = 1 2    

3 6 c1 = − c0 = − c0  5 5 2   5 5 6 6 c2 = − c 1 = − − c0 = c0  7 7 5 7   7 14 6 4 c = − c0  c3 = − c2 = − 27 27 7 0 9 2 and so on. One Frobenius solution is

  6 4 6 y1
x = c0 x − x2 + x3 − x4 + · · ·  5 7 9

Because r1 is a nonnegative integer, this first Frobenius solution is actually a power series about 0. 1 For a second Frobenius solution, substitute r = r2 = − into the recurrence relation. To 2 avoid confusion we will replace cn with cn∗ in this relation. We get   3 1+2 n− 2     cn∗ = −  c∗ 1 3 1 1 1 n−1 n− n− + n− − 2 2 2 2 2 for n = 1 2     This simplifies to 2n − 2 ∗ c  cn∗ = −  3 n−1 n n− 2 It happens in this example that c1∗ = 0, so each cn∗ = 0 for n = 1 2    and the second Frobenius solution is y2
x =



cn∗ xn−1/2 = c0∗ x−1/2

for x > 0


n=0

The method of Frobenius is justified by the following theorem.

170

CHAPTER 4

Series Solutions

THEOREM 4.3

Method of Frobenius

Suppose x0 is a regular singular point of P
xy

+ Q
xy
+ R
xy = 0. Then there exists at least one Frobenius solution  y
x = cn
x − x0 r n=0

with c0 = 0. Further, if the Taylor expansions of
x − x0 Q
x/R
x and
x − x0 2 R
x/P
x about x0 converge in an open interval
x0 − h x0 + h, then this Frobenius series also converges in this interval, except perhaps at x0 itself. It is significant that the theorem only guarantees the existence of one Frobenius solution. Although we obtained two such solutions in the preceding example, the next example shows that there may be only one.

EXAMPLE 4.11

Suppose we want to solve x2 y

+ 5xy
+
x + 4y = 0 Zero is a regular singular point, so attempt a Frobenius solution y
x = into the differential equation to get 


n + r
n + r − 1cn xn+r +

n=0



5
n + rcn xn+r +

n=0





cn xn+r+1 +

n=0

n=0 cn x



n+r

. Substitute

4cn xn+r = 0

n=0

Shift indices in the third summation to write this equation as 


n + r
n + r − 1cn xn+r +

n=0



5
n + rcn xn+r +

n=0



cn−1 xn+r +

n=1



4cn xn+r = 0

n=0

Now combine terms to write

r
r − 1 + 5r + 4c0 xr +




n + r
n + r − 1cn + 5
n + rcn + cn−1 + 4cn xn+r = 0

n=1

Setting the coefficient of xr equal to zero (since c0 = 0 as part of the method), we get the indicial equation r
r − 1 + 5r + 4 = 0 with the repeated root r = −2. The coefficient of xn+r in the series, with r = −2 inserted, gives us the recurrence relation
n − 2
n − 3cn + 5
n − 2cn + cn−1 + 4cn = 0 or cn = −

1 c
n − 2
n − 3 + 5
n − 2 + 4 n−1

for n = 1 2     This simplifies to cn = −

1 c n2 n−1

for n = 1 2 3    

4.3 Singular Points and the Method of Frobenius

171

Some of the coefficients are c1 = −c0 1 1 1 c2 = − c 1 = c 0 = c 4 4
22 0 1 1 1 c =− c3 = − c 2 = − c 9 4·9 0
2 · 32 0 c4 = −

1 1 1 c c3 = c0 = 16 4 · 9 · 16
2 · 3 · 42 0

and so on. In general, cn =
−1n

1 c
n!2 0

for n = 1 2 3     The Frobenius solution we have found is y
x = c0 x−2 − x−1 + = c0

 n=0


−1n

1 1 2 1 − x+ x +··· 4 36 576

1 n−2 x 
n!2

In this example, xQ
x/P
x = x
5x/x2  = 5 and x2 R
x/P
x = x2
x + 4/x2 = x + 4. These polynomials are their own Maclaurin series about 0, and these series, being finite, converge for all x. By Theorem 4.3, the Frobenius series solution converges for all x, except x = 0. In this example the method of Frobenius produces only one solution. In the last example the recurrence relation produced a simple formula for cn in terms of c0 . Depending on the coefficients in the differential equation, a formula for cn in terms of c0 may be quite complicated, or it may even not be possible to write a formula in terms of elementary algebraic expressions. We will give another example, having some importance for later work, in which the Frobenius method may produce only one solution.

EXAMPLE 4.12 Bessel Functions of the First Kind

The differential equation x2 y

+ xy
+
x2 −  2 y = 0 is called Bessel’s equation of order , for  ≥ 0. Although it is a second-order differential equation, this description of it as being of order  refers to the parameter  appearing in it, and is traditional. Solutions of Bessel’s equation are called Bessel functions, and we will encounter them in Chapter 16 when we treat special functions, and again in Chapter 18 when we analyze heat conduction in an infinite cylinder. Zero is a regular singular point of Bessel’s equation, so attempt a solution y
x =

 n=0

cn xn+r 

172

CHAPTER 4

Series Solutions

Upon substituting this series into Bessel’s equation, we obtain

r
r − 1 + r −  2 c0 xr + r
r + 1 +
r + 1 −  2 c1 xr+1 +




n + r
n + r − 1 +
n + r −  2 cn + cn−2 xn+r = 0

(4.15)

n=2

Set the coefficient of each power of x equal to zero. Assuming that c0 = 0, we obtain the indicial equation r 2 −  2 = 0 with roots ±. Let r =  in the coefficient of xr+1 in equation (4.15) to get
2 + 1c1 = 0 Since 2 + 1 = 0, we conclude that c1 = 0. From the coefficient of xn+r in equation (4.15), we get


n + r
n + r − 1 +
n + r −  2 cn + cn−2 = 0 for n = 2 3     Set r =  in this equation and solve for cn to get cn = −

1 c n
n + 2 n−2

for n = 2 3     Since c1 = 0, this equation yields c3 = c5 =    = codd = 0 For the even-indexed coefficients, write c2n = − =− =

1 1 c2n−2 = − 2 c 2n
2n + 2 2 n
n +  2n−2 1 −1 c 22 n
n +  2
n − 1
2
n − 1 + 2 2n−4 1

c 24 n
n − 1
n + 
n +  − 1 2n−4

=  = =


−1n c 22n n
n − 1 · · ·
2
1
n + 
n − 1 +  · · ·
1 +  0
−1n

c 22n n!
1 + 
2 +  · · ·
n +  0

One Frobenius solution of Bessel’s equation of order  is therefore y1
x = c0




−1n x2n+  2n n=0 2 n!
1 + 
2 +  · · ·
n + 

(4.16)

These functions are called Bessel functions of the first kind of order . The roots of the indicial equation for Bessel’s equation are ±. Depending on , we may or may not obtain two linearly independent solutions by using  and − in the series solution (4.16). We will discuss this in more detail when we treat Bessel functions in Chapter 16, where we will see that, when  is a positive integer, the functions obtained by using  and then − in the recurrence relation are linearly dependent.

4.4 Second Solutions and Logarithm Factors

PROBLEMS

SECTION 4.3

In each of Problems 1 through 6, find all of the singular points and classify each singular point as regular or singular.

rence relation, and (d) use the results of (b) and (c) to find the first five nonzero terms of two linearly independent Frobenius solutions. 7. 4x2 y

+ 2xy
− xy = 0

1. x2
x − 32 y

+ 4x
x2 − x − 6y
+
x2 − x − 2y = 0 2.
x3 − 2x2 − 7x − 4y

− 2
x2 + 1y
+
5x2 − 2xy = 0 3. x2
x − 2y

+
5x − 7y
+ 2
3 + 5x2 y = 0

8. 16x2 y

− 4x2 y
+ 3y = 0 9. 9x2 y

+ 2
2x + 1y = 0

4.
9 − x2 y

+
2 + x2 y = 0

10. 12x2 y

+ 5xy
+
1 − 2x2 y = 0

5.
x − 2−1 y

+ x−5/2 y = 0

11. 2xy

+
2x + 1y
+ 2y = 0

6. x2 sin2
x − y

+ tan
x −  tan
xy
+
7x − 2 cos
xy = 0

12. 2x2 y

− xy
+
1 − x2 y = 0

In each of Problems 7 through 15, (a) show that zero is a regular singular point of the differential equation, (b) find and solve the indicial equation, (c) determine the recur-

14. 3x2 y

+ 4xy
−
3x + 2y = 0

4.4

173

13. 2x2 y

+ x
2x + 1y
−
2x2 + 1y = 0 15. 9x2 y

+ 9xy
+
9x2 − 4y = 0

Second Solutions and Logarithm Factors In the preceding section we saw that under certain conditions we can always produce a Frobenius series solution of equation (4.10), but possibly not a second, linearly independent solution. This may occur if the indicial equation has a repeated root, or even if it has distinct roots that differ by a positive integer. In the case that the method of Frobenius only produces one solution, there is a method for finding a second, linearly independent solution. The key is to know what form to expect this solution to have, so that this template can be substituted into the differential equation to determine the coefficients. This template is provided by the following theorem. We will state the theorem with x0 = 0 to simplify the notation. To apply it to a differential equation having a singular point x0 = 0, use the change of variables z = x − x0 .

THEOREM 4.4

A Second Solution in the Method of Frobenius

Suppose 0 is a regular singular point of P
xy

+ Q
xy
+ R
xy = 0 Let r1 and r2 be roots of the indicial equation. If these are real, suppose r1 ≥ r2 . Then 1. If r1 − r2 is not an integer, there are two linearly independent Frobenius solutions y1
x =

 n=0

cn xn+r1

and

y2
x =



cn∗ xn+r2 

n=0

with c0 = 0 and c0∗ = 0. These solutions are valid in some interval
0 h or
−h 0.

174

CHAPTER 4

Series Solutions

2. If r1 − r2 = 0, there is a Frobenius solution y1
x = a second solution y2
x = y1
x ln
x +





n=0 cn x

n+r1

with c0 = 0 as well as

cn∗ xn+r1 

n=1

Further, y1 and y2 form a fundamental set of solutions on some interval
0 h. 3. If r1 − r2 is a positive integer, then there is a Frobenius series solution y1
x =



cn xn+r1 

n=0

In this case there is a second solution of the form y2
x = ky1
x ln
x +



cn∗ xn+r2 

n=0

If k = 0 this is a second Frobenius series solution; if not the solution contains a logarithm term. In either event, y1 and y2 form a fundamental set on some interval
0 h. We may now summarize the method of Frobenius as follows, for the equation P
xy

+ 0 is a regular singular point. Q
xy
+ R
xy = 0. Suppose  n+r c into the differential equation. From the indicial equation, Substitute y
x = n=0 n x determine the values of r. If these are distinct and do not differ by an integer, we are guaranteed two linearly independent Frobenius solutions. If the indicial equation has repeated roots, then there is just one Frobenius solution y1 . But there is a second solution y2
x = y1
x ln
x +



cn∗ xn+r1 

n=1

The series on the right starts its summation at n = 1, not n = 0. Substitute y2
x into the differential equation and obtain a recurrence relation for the coefficients cn∗ . Because this solution has a logarithm term, y1 and y2 are linearly independent. If r1 − r2 is a positive integer, there may or may not be a second Frobenius solution. In this case there is a second solution of the form y2
x = ky1
x ln
x +



cn∗ xn+r2 

n=0

Substitute y2 into the differential equation and obtain an equation for k and a recurrence relation for the coefficients cn∗ . If k = 0, we obtain a second Frobenius solution; if not, then this second solution has a logarithm term. In either case y1 and y2 are linearly independent. In the preceding section we saw in Example 4.10 a differential equation in which r1 − r2 was not an integer. There we found two linearly independent Frobenius solutions. We will illustrate cases (2) and (3) of the theorem.

EXAMPLE 4.13 Conclusion (2), Equal Roots

Consider again x2 y

+ 5xy
+
x + 4y = 0. In Example 4.11 we found one Frobenius solution y1
x = c0

 n=0


−1n

1 n−2 x 
n!2

4.4 Second Solutions and Logarithm Factors

175

The indicial equation is
r + 22 = 0 with the repeated root r = −2. Conclusion (2) of the theorem suggests that we attempt a second solution of the form y2
x = y1
x ln
x +



cn∗ xn−2 

n=1

Note that the series on the right begins at n = 1, not n = 0. Substitute this series into the differential equation to get, after some rearrangement of terms, 4y1 + 2xy1
+




n − 2
n − 3cn∗ xn−2 +

n=1

+





5
n − 2cn∗ xn−2 +

n=1



cn∗ xn−1

n=1

4cn∗ xn−2 + ln
x x2 y1

+ 5xy1
+
x + 4y1  = 0

n=1

The bracketed coefficient of ln
x is zero because y1 is a solution of the differential equation. In the last equation, c0∗ = 1 (we need only one second solution), shift indices to write  choose  ∗ n−1 ∗ n−2 c x = c x , and substitute the series obtained for y1
x to get n=1 n n=2 n−1  4
−1n 2
−1n −1 ∗ −1 +
n − 2 − 2x + c1 x +
n!2
n!2 n=2 ∗ + 4cn∗  xn−2 = 0 +
n − 2
n − 3cn∗ + 5
n − 2cn∗ + cn−1

Set the coefficient of each power of x equal to zero. From the coefficient of x−1 we get c1∗ = 2 From the coefficient of xn−2 in the summation we get, after some routine algebra, 2
−1n ∗ n + n2 cn∗ + cn−1 = 0
n!2 or cn∗ = −

1 ∗ 2
−1n c − n2 n−1 n
n!2

for n = 2 3 4     This enables us to calculate as many coefficients as we wish. Some of the terms of the resulting solution are y2
x = y1
x ln
x +

11 25 2 2 3 137 − + x− x + x3 + · · ·  x 4 108 3456 432,000

Because of the logarithm term, it is obvious that this solution is not a constant multiple of y1 , so y1 and y2 form a fundamental set of solutions (on some interval
0 h). The general solution is y
x = C1 + C2 ln
x + C2


−1n n=0


n!2

xn−2

2 3 11 25 2 137 − + x− x + x3 + · · ·  x 4 108 3456 432,000

176

CHAPTER 4

Series Solutions

EXAMPLE 4.14 Conclusion (3), with k = 0

 n+r The equation x2 y

+x2 y
−2y = 0 has a regular singular point at 0. Substitute y
x = n=0 cn x and shift indices to obtain 

r
r − 1 − 2c0 xr +
n + r
n + r − 1cn +
n + r − 1cn−1 − 2cn xn+r = 0 n=1

Assume that c0 = 0. The indicial equation is r 2 − r − 2 = 0, with roots r1 = 2 and r2 = −1. Now r1 − r2 = 3 and case (3) of the theorem applies. For a first solution, set the coefficient of xn+r equal to zero to get
n + r
n + r − 1cn +
n + r − 1cn−1 − 2cn = 0

(4.17)

Let r = 2 to get
n + 2
n + 1cn +
n + 1cn−1 − 2cn = 0 or cn = −

n+1 c n
n + 3 n−1

for n = 1 2    

Using this recurrence relation to generate terms of the series, we obtain

3 2 1 3 1 4 1 5 1 6 1 2 x − x + x +···  y1
x = c0 x 1 − x + x − x + 2 20 30 168 1120 8640 Now try the second root r = −1 in the recurrence relation (4.17). We get ∗ − 2cn∗ = 0
n − 1
n − 2cn∗ +
n − 2cn−1

for n = 1 2     When n = 3, this gives c2∗ = 0, which forces cn∗ = 0 for n ≥ 2. But then 1 y2
x = c0∗ + c1∗  x Substitute this into the differential equation to get   2 ∗ −3 2 ∗ −2 ∗ ∗1 x
2c0 x  + x
−c0 x  − 2 c1 + c0 = −c0∗ − 2c1∗ = 0 x Then c1∗ = − 21 c0∗ and we obtain the second solution   1 1 − y2
x = c0∗  x 2 with c0∗ nonzero but otherwise arbitrary. The functions y1 and y2 form a fundamental set of solutions.

EXAMPLE 4.15 Conclusion (3), k = 0

Consider differential equation xy

− y = 0, which has a regular singular point at 0. Substitute the y
x = n=0 cn xn+r to obtain  n=0


n + r
n + r − 1cn xn+r−1 −

 n=0

cn xn+r = 0

4.4 Second Solutions and Logarithm Factors

177

Shift indices in the second summation to write this equation as
r 2 − rc0 xr−1 +




n + r
n + r − 1cn − cn−1 xn+r−1 = 0

n=1

The indicial equation is r 2 − r = 0, with roots r1 = 1, r2 = 0. Here r1 − r2 = 1, a positive integer, so we are in case (3) of the theorem. The recurrence relation is
n + r
n + r − 1cn − cn−1 = 0 for n = 1 2     Let r = 1 and solve for cn  cn =

1 c n
n + 1 n−1

for n = 1 2 3    

Some of the coefficients are c1 =

1 c 1
2 0

c2 =

1 1 c = c 2
3 1 2
2
3 0

c3 =

1 1 c2 = c 3
4 2
3
2
3
4 0

In general, we find that cn =

1 c n!
n + 1! 0

for n = 1 2     This gives us a Frobenius series solution 

1 xn+1 n!
n + 1! n=0

1 1 1 4  x +  = c0 x + x 2 + x 3 + 2 12 144

y1
x = c0

In this example, if we put r = 0 into the recurrence relation, we get n
n − 1cn − cn−1 = 0 for n = 1 2     But if we put n = 1 into this equation, we get c0 = 0, contrary to the assumption that c0 = 0. Unlike the preceding example, we cannot find a second Frobenius solution by simply putting r2 into the recurrence relation. Try a second solution y2
x = ky1
x ln
x +



cn∗ xn

n=0

(here xn+r2 = xn because r2 = 0). Substitute this into the differential equation to get  x

ky1



1  − ky1 2 + n
n − 1cn∗ xn−2 x x n=2

1 ln
x + 2ky1


− ky1 ln
x −

 n=0

cn∗ xn = 0



(4.18)

178

CHAPTER 4

Series Solutions

Now k ln
x xy1

− y1  = 0 because y1 is a solution of the differential equation. For the remaining terms in equation (4.18), insert the series for y1
x (with c0 = 1 for convenience) to get 2k



   1 n 1 n ∗ n−1 x x − k + c n
n − 1x − cn∗ xn = 0 n 2 n=0
n! n=0 n!
n + 1! n=2 n=0

Shift indices in the third summation to write this equation as 2k



   1 n 1 n ∗ n x x − k + c
n + 1nx − cn∗ xn = 0 n+1 2 n=0
n! n=0 n!
n + 1! n=1 n=0

Then
2k − k − c0∗ x0 +

 n=1



2k k ∗ ∗ + n
n + 1cn+1 − cn xn = 0 −
n!2 n!
n + 1!

Then k − c0∗ = 0 and, for n = 1 2     2k k ∗ + n
n + 1cn+1 − − cn∗ = 0
n!2 n!
n + 1! This gives us k = c0∗ and the recurrence relation

1
2n + 1k ∗ cn+1 = cn∗ − n
n + 1 n!
n + 1! for n = 1 2 3     Since c0∗ can be any nonzero real number, we may choose c0∗ = 1. Then k = 1 For a particular second solution, let c1∗ = 0 obtaining: 3 7 35 4 x −···  y2
x = y1
x ln
x + 1 − x2 − x3 − 4 36 1728 To conclude this section, we will produce a second solution for Bessel’s equation, in a case where the Frobenius method yields only one solution. This will be of use later when we study Bessel functions.

EXAMPLE 4.16 Bessel Function of the Second Kind

Consider Bessel’s equation of zero order
 = 0. From Example 4.12, this is x2 y

+ xy
+ x2 y = 0 We know from that example that the indicial equation has only one root, r = 0. From equation (4.16), with c0 = 1, one Frobenius solution is y1
x =

 k=0


−1k

1 22k
k!2

x2k 

Attempt a second, linearly independent solution of the form y2
x = y1
x ln
x +

 k=1

ck∗ xk 

4.4 Second Solutions and Logarithm Factors

179

Substitute y2
x into the differential equation to get  1 xy1

ln
x + 2y1
− y1 + k
k − 1ck∗ xk−1 x k=2   1 + y1
ln
x + y1 + kck∗ xk−1 + xy1 ln
x + ck∗ xk+1 = 0 x k=1 k=1

Terms involving ln
x and y1
x cancel, and we are left with 2y1
+



k
k − 1ck∗ xk−1 +

k=2



kck∗ xk−1 +

k=1



ck∗ xk+1 = 0

k=1

Since k
k − 1 = k − k, part of the first summation cancels all terms except the k = 1 term in the second summation, and we have 2

2y1
+



k2 ck∗ xk−1 + c1∗ +

k=2



ck∗ xk+1 = 0

k=1

Substitute the series for y1
into this equation to get 2



 
−1k 2k−1 2 ∗ k−1 ∗ x + k c x + c + ck∗ xk+1 = 0 k 1 2k−1 k!
k − 1! k=1 2 k=2 k=1

Shift indices in the last series to write this equation as 


−1k

k=1

22k−2 k!
k − 1!

x2k−1 + c1∗ + 4c2∗ x +



∗
k2 ck∗ + ck−2 xk−1 = 0

(4.19)

k=3

The only constant term on the left side of this equation is c1∗ , which must therefore be zero. The only even powers of x appearing in equation (4.19) are in the right-most series when k is odd. The coefficients of these powers of x must be zero, hence ∗ k2 ck∗ + ck−2 =0

for k = 3 5 7    

But then all odd-indexed coefficients are multiples of c1∗ , which is zero, so ∗ c2k+1 =0

for k = 0 1 2    

To determine the even-indexed coefficients, replace k by 2j in the second summation of equation (4.19) and k with j in the first summation to get 


−1j 2j−1 ∗ ∗ ∗ x + 4c x +
4j 2 c2j + c2j−2 x2j−1 = 0 2 2j−2 j!
j − 1! j=1 2 j=2

Now combine terms and write this equation as


−1j 2 ∗ ∗ 2j−1
4c2∗ − 1x + + 4j c + c = 0 2j 2j−2 x 2j−2 j!
j − 1! 2 j=2 Equate the coefficient of each power of x to zero. We get c2∗ =

1 4

and the recurrence relation ∗ = c2j


−1j+1 1 ∗ c − 22j j!2 j 4j 2 2j−2

for j = 2 3 4    

CHAPTER 4

180

Series Solutions

If we write some of these coefficients, a pattern emerges:

1 −1 ∗  c4 = 2 2 1 + 24 2

1 1 1 ∗ c6 = 2 2 2 1 + +  246 2 3 and, in general, ∗ = c2j



1
−1j+1 1
−1j+1 + · · · + ∅
j 1 + = 22 42 · · ·
2j2 2 j 22j
j!2

where ∅
j = 1 +

1 1 +···+ 2 j

for j = 1 2 · · · 

We therefore have a second solution of Bessel’s equation of order zero: y2
x = y1
x ln
x +


−1k+1 k=1

22k
k!2

∅
kx2k

for x > 0. This solution is linearly independent from y1
x for x > 0. When a differential equation with a regular singular point has only one Frobenius series solution expanded about that point, it is tempting to try reduction of order to find a second solution. This is a workable strategy if we can write y1
x in closed form. But if y1
x is an infinite series, it may be better to substitute the appropriate form of the second solution from Theorem 4.4 and solve for the coefficients.

SECTION 4.4

PROBLEMS

In each of Problems 1 through 10, (a) find the indicial equation, (b) determine the appropriate form of each of two linearly independent solutions, and (c) find the first five terms of each of two linearly independent solutions. In Problems 11 through 16, find only the form that two linearly independent solutions should take. 1. xy

+
1 − xy
+ y = 0 2. xy

− 2xy
+ 2y = 0

9. x
2 − xy

− 2
x − 1y
+ 2y = 0 10. x2 y

+ x
x3 + 1y
− y = 0

  4 y=0 11. 25x
1 − x2 y

− 20
5x − 2y
+ 25x − x   16


y + 12. 6
3x − 4
5x + 8y +
2x − 21 3 + x   27 4− 2 y = 0 x

6. 4x2 y

+ 4xy
− y = 0




13. 12x
4  + 3xy − 2
5x + 7
7x − 2y + 1 24 5 − y=0 3x   8 y=0 14. 3x
2x + 3y

+ 2
6 − 5xy
+ 7 2x − x

7. x2 y

− 2xy
−
x2 − 2y = 0

15. x
x + 4y

− 3
x − 2y
+ 2y = 0

3. x
x − 1y

+ 3y
− 2y = 0 4. 4x2 y

+ 4xy
+
4x2 − 9y = 0 5. 4xy

+ 2y
+ y = 0








8. xy − y + 2y = 0

16.
3x3 + x2 y

− x
10x + 1y
+
x2 + 2y = 0

CHAPTER

5

EULER’S METHOD ONE-STEP METHODS MULTISTEP METHODS EULER’S METHOD ONE-STEP METHODS MULTISTEP METHODS EULER’S METHOD ONE-STEP METHODS MULTISTEP METHODS EULER’S METHOD

Numerical Approximation of Solutions

Often we are unable to produce a solution of an initial value problem in a form suitable for drawing a graph or calculating numerical values. When this happens we may turn to a scheme for approximating numerical values of the solution. Although the idea of a numerical approximation is not new, it is the development and ready accessibility of high speed computers that have made it the success that it is today. Some problems thought to be intractable thirty years ago are now considered solved from a practical point of view. Using computers and numerical approximation techniques, we now have increasingly accurate models for weather patterns, national and international economies, global warming, ecological systems, fluid flow around airplane wings and ship hulls, and many other phenomena of interest and importance. A good numerical approximation scheme usually includes the following features. 1. At least for first order initial value problems, the scheme usually starts at a point x0 where the initial value is prescribed, then builds approximate values of the solution at points specified to the left or right of x0 . The accuracy of the method will depend on the distance between successive points at which the approximations are made, their increasing distance from x0 , and of course the coefficients in the differential equation. Accuracy can also be influenced by the programming and by the architecture of the computer. For some complex models, such as the Navier-Stokes equations governing fluid flow, computers have been built with architecture dedicated to efficient approximation of solutions of that particular model. 2. A good numerical scheme includes an estimate or bound on the error in the approximation. This is used to understand the accuracy in the approximation, and often to guide the user in choosing certain parameters (such as the number of points at which approximations are made, and the distance between successive points). Often a compromise must be made between increasing accuracy (say, by choosing more points) and keeping the time or cost of the computation within reason. The type of problem under consideration may dictate what might be acceptable bounds on the error. If NASA is placing a satellite in a Jupiter orbit, a one meter error might be acceptable, while an error of this magnitude would be catastrophic in performing eye surgery. 181

182

CHAPTER 5

Numerical Approximation of Solutions

3. The method must be implemented on a computer. Only simple examples, devised for illustrative purposes, can be done by hand. Many commercially available software packages include routines for approximating and graphing solutions of differential equations. Among these are MAPLE, MATHEMATICA and MATLAB. We will now develop some specific methods.

Euler’s Method Euler’s method is a scheme for approximating the solution of y
= f
x y

y
x0  = y0 

in which x0 , y0 and the function f are given. The method is a good introduction to numerical schemes because it is conceptually simple and geometrically appealing, although it is not the most accurate. Let y
x denote the solution (which we know exists, but do not know explicitly). The key to Euler’s method is that if we know y
x at some x, then we can compute f
x y
x, and therefore know the slope y

x of the tangent to the graph of the solution at that point. We will exploit this fact to approximate solution values at points x1 = x0 + h, x2 = x0 + 2h,    , xn = x0 + nh. First choose h (the step size) and the number n of iterations to be performed. Now form the first approximation. We know y
x0  = y0 . Calculate f
x0  y0  and draw the line having this slope through
x0  y0 . This line is tangent to the integral curve through
x0  y0 . Move along this tangent line to the point
x1  y1 . Use y1 as an approximation to y
x1 . This is illustrated in Figure 5.1. We have some hope that this is a "good" approximation, for h "small", because the tangent line fits the curve closely "near" the point. Next compute f
x1  y1 . This is the slope of the tangent to the graph of the solution of the differential equation passing through
x1  y1 . Draw the line through
x1  y1  having this slope, and move along this line to
x2  y2 . This determines y2 , which we take as an approximation to y
x2  (see Figure 5.1 again). Continue in this way. Compute f
x2  y2  and draw the line with this slope through
x2  y2 . Move along this line to
x3  y3  and use y3 as an approximation to y
x3 . pe slo ) ,y2 x2 f(

y

(x2, y2)

(x1, y1)

slope y) f(x 1, 1

(x3, y3) (x0, y0)

x1 x2

x x3

slo pe 0, y 0)

x0

f (x

5.1

FIGURE 5.1 Approximation points formed according to Euler’s method.

5.1 Euler’s Method

183

In general, once we have reached
xk  yk , draw the line through this point having slope f
xk  yk  and move along this line to
xk+1  yk+1 . Take yk+1 as an approximation to y
xk+1 . This is the idea of the method. Obviously it is quite sensitive to how much f
x y changes if x and y are varied by a small amount. The method also tends to accumulate error, since we use the approximation yk to make the approximation yk+1 . In Figure 5.2, the successively drawn line segments used to determine the approximate values move away from the actual solution curve as x increases, causing the approximations to be less accurate as more of them are made (that is, as n is chosen larger). Following segments of lines is conceptually simple and appealing, but it is not sophisticated enough to be very accurate in general. We will now derive an analytic expression for the approximate solution value yk at xk . From Figure 5.1, y1 = y0 + f
x0  y0 
x1 − x0  At the next step, y2 = y1 + f
x1  y1 
x2 − x1  After we have obtained the approximate value yk , the next step (Figure 5.3) gives yk+1 = yk + f
xk  yk 
xk+1 − xk  Since each xk+1 − xk = h, we can summarize the discussion as follows.

DEFINITION 5.1

Euler’s Method

Euler’s method is to define yk+1 in terms of yk by yk+1 = yk + f
xk  yk 
xk+1 − xk  or yk+1 = yk + hf
xk  yk  for k = 0 1 2     n − 1. yk is the Euler approximation to y
xk .

y (x4, y4) y (x3, y3)

yk+1

(x2, y2) (x0, y0)

x0

slope f (xk, yk) yk

(x1, y1)

x1

x2

x3

x4

x

FIGURE 5.2 Accumulating error in Euler’s method.

xk xk+1 FIGURE 5.3

x

184

CHAPTER 5

Numerical Approximation of Solutions

EXAMPLE 5.1

Consider √ y
= x y

y
2 = 4

This separable differential equation is easily solved: 2  x2 y
x = 1 +  4 This enables us to observe how the method works by direct comparison with the exact solution. First we must decide on h and n. Since we do not have any error estimates, we have no rationale for making a particular choice. For illustration, choose h = 02 and n = 20. Then x0 = 2 and x20 = x0 + nh = 2 +
20
02 = 6. Now √ yk+1 = yk + 02xk yk for k = 0 1 2     19 Table 5.1 lists the Euler approximations, and Figure 5.4 shows a graph of this approximate solution (actually, a smooth curve drawn through the approximated points), together with a graph of the actual solution, for comparison. Notice that the approximation becomes less accurate as x moves further from 2.

TA B L E 5 . 1

√ Approximate Values of the Solution of y
= x yy
2 = 4h = 02n = 20

x

yapp
x

x

yapp
x

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

4 48 5763991701 6916390802 828390462 9895723242 11783171355 1398007530 1652259114 1944924644 2280094522

4.2 4.4 4.6 4.8 5.0 5.2 5.4 5.6 5.8 6.0

26.62097204 30.95499533 35.85107012 41.35964033 47.53354060 54.42799784 62.10063249 70.61145958 80.02288959 90.39972950

y(1+

100

x2 4

2

)

80 60

Euler approximation

40 20 0 2

3

4

5

6

x

FIGURE 5.4 Exact and Euler approximate solutions of √ y
= x y y
2 = 4 with stepsize 0.2 and twenty iterations.

5.1 Euler’s Method

TA B L E 5 . 2

185

√ Approximate Values of the Solution of y
= x yy
2 = 4h = 01n = 40

x

yapp
x

x

yapp
x

2.0 2.1 2.2 2.3 2.4 2.5 2.6 2.7 2.8 2.9 3.0 3.1 3.2 3.3 3.4 3.5 3.6 3.7 3.8 3.9 4.0

4 44 44840499716 5324524701 5855248129 6435991009 7070222385 7761559519 8513768060 9330762195 1021660479 1117550751 1221183095 1333008472 1453492754 1583116734 1722376132 1871781603 2031858741 2203148088 2386205135

4.1 4.2 4.3 4.4 4.5 4.6 4.7 4.8 4.9 5.0 5.1 5.2 5.3 5.4 5.5 5.6 5.7 5.8 5.9 6.0

25.8160330 27.89919080 30.11761755 32.47743693 34.98495199 37.64661553 40.46903007 43.45894792 46.62327117 49.96905171 53.50349126 57.23394138 61.16790347 65.31302881 69.67711854 74.26812370 79.09414521 84.16343390 89.48439050 95.06556568

Exact solution

100

h = 0.1 h = 0.2

80 60 40 20 0 2

3

4

5

6

x

FIGURE 5.5 Exact and Euler approximate solutions of √ y
= x y y
2 = 4, first with h = 02 and twenty iterations, and then h = 01 and forty iterations.

The accuracy of this method depends on h. If we choose h = 01 and n = 40 (so the approximation is still for 2 ≤ x ≤ 6), we get the approximate values of Table 5.2. A graph of this approximation is shown in Figure 5.5, showing an improved approximation by choosing h smaller. With today’s computing power, we would have no difficulty using a much smaller h.

186

CHAPTER 5

Numerical Approximation of Solutions

EXAMPLE 5.2

Consider y
= sin
xy

y
2 = 1

We cannot write a simple solution for this problem. Figure 5.6 shows a direction field for y
= sin
xy, and Figure 5.7 repeats this direction field, with some integral curves, including the one through
2 1. This is a graph of the solution (actually an approximation done by the software used for the direction field). For a numerical approximation of the solution, choose h = 02 and n = 20 to obtain an approximate solution for 2 ≤ x ≤ 6. The generated values are given in Table 5.3, and a smooth curve is drawn through these points in Figure 5.8. y 4

2

−2

−4

0

4

2

x

–2

–4

A direction field for y
= sin
xy.

FIGURE 5.6

y 6 5 4 3 2 (2, 1) 1

0

1

2

3

4

5

FIGURE 5.7 A direction field and some integral curves for y
= sin
xy, including the integral curve through (2,1).

6

x

5.1 Euler’s Method

TA B L E 5 . 3

187

Approximate Values of the Solution of y
=sin
xyy
2=1h=02n=20

x

yapp
x

x

yapp
x

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 2.6 3.8 4.0

1 1.181859485 1.284944186 1.296483096 1.251031045 1.180333996 1.102559149 1.027151468 0.9584362156 0.8976573370 0.8444064257

4.2 4.4 4.6 4.8 5.0 5.2 5.4 5.6 5.8 6.0

0.7976369224 0.7562418383 0.7192812325 0.6860163429 0.6558744703 0.6248056314 0.6032491169 0.5801105326 0.5587461082 0.5389516129

y 1.2 1.0 0.8 0.6 x 0

1

2

3

4

5

6

FIGURE 5.8 Euler approximate solution of y
= sin
xy y
2 = 1
h = 02 n = 20.

From the examples, it appears that the error in an Euler approximation is proportional to h. It can be shown that this is indeed the case, and for this reason Euler’s method is a first order method. If a method has error that is proportional to hp , it is called an order p method.

5.1.1

A Problem in Radioactive Waste Disposal

Disposal of radioactive waste materials generated by nuclear power plants, medical research, military testing and other sources is a serious international problem. In view of the long halflives of some of the materials involved, there is no real prospect for disposal, and the problem becomes one of safe storage. For example, Uranium-235 has a half-life of 713
108  years, and Thorium-232 a half-life of 139
1010  years. If a ton of Thorium-232 is stored, more than half of it will still be here to see our sun consume all of its fuel and die. Storage plans have been proposed and studied on both national levels (for example, by the U.S. Atomic Energy Commission) and the international level (by the International Atomic Energy Agency of the United Nations). Countries have developed a variety of policies and plans. Argentina has embarked on a program of storing containers of radioactive materials in granite vaults. Belgium is planning to bury containers in clay deposits. Canada plans to use crystalline rocks in the Canadian shield. The Netherlands is considering salt domes. France and Japan are planning undersea depositories. And the United States has a diversified approach

188

CHAPTER 5

Numerical Approximation of Solutions

which have included sites at Yucca Mountain in Nevada and the Hanford Reservation in the state of Washington. One idea which has been considered is to store the material in fifty-five gallon containers and drop them into the ocean at a point about 300 feet deep, shallow enough to prevent rupture of the drums by water pressure. It was found that drums could be manufactured that would endure indefinitely at this depth. But then another point was raised. Would the drums withstand the impact of settling on the ocean floor after being dropped from a ship? Testing showed that the drums could indeed rupture if they impacted the bottom at a speed in excess of 40 feet per second. The question now is: will a drum achieve this velocity in a 300 foot descent through seawater? To answer this question, we must analyze what happens when a drum is dropped into the water and allowed to settle to the bottom. Each 55-gallon drum weighs about 535 pounds after being filled with the material and some insulation. When in the water, the drum is buoyed up by a force equal to the weight of the water displaced. Fifty-five gallons is about 7.35 cubic feet, and the density of seawater is about 64 pounds per cubic foot, so each barrel will be subject to a buoyant force of about 470 pounds. In addition to this buoyant force, the water will impose a drag on the barrel as it sinks, impeding its descent. It is well known that objects sinking in a fluid are subject to a drag force which is proportional to a power of the velocity. Engineers had to determine the constant of proportionality and the exponent for a drum in seawater. √ After testing, they estimated that the drag force of the water was approximately equal to 05v 10/3 pounds, in which v is the velocity in feet per second. Let y
t be the depth of the drum in the water at time t, with downward chosen as the positive direction. Let y = 0 at the (calm) surface of the water. Then v
t = y

t. The forces acting on the drum are the buoyant and drag forces (acting upward) and the force of gravity (acting downward). Since the force of gravity has magnitude mg, with m the mass of the drum, then by Newton’s law, m

√ dv = mg − 470 − 05v 10/3  dt

For this problem, mg = 535 pounds. Use g = 32 ft/sec2 to determine that m = 167 slugs. Assume that the drum is released from rest at the surface of the water. The initial value problem for the velocity of the descending drum is 167

√ dv = 535 − 470 − 05v 10/3  dt

v
0 = 0

or  √ dv 1  = 65 − 05v 10/3  dt 167

v
0 = 0

We want the velocity with which the drum hits bottom. One approach might give us a quick answer. It is not difficult to show that a drum sinking in seawater will have a terminal velocity. If the terminal velocity of the drum is less than 40 feet per second, then a drum released from rest will never reach a speed great enough to break it open upon impact with the ocean floor, regardless of the depth! Unfortunately, a quick calculation, letting dv/dt = 0, shows that the terminal velocity is about 100 feet per second, not even close to 40. This estimate is therefore inconclusive in determining whether the drums have a velocity of 40 feet per second upon impact at 300 feet. We could try solving the differential equation for v
t and integrating to get an equation for the depth at time t. Setting y
t = 300 would then yield the time required for the drum to reach this depth, and we could put this time back into v
t to see if the velocity exceeds

5.1 Euler’s Method

189

40 feet per second at this time. The differential equation is separable. However, solving it leads to the integral
1 √ dv 10/3 v − 130 which has no elementary evaluation. Another approach would be to express the velocity as a function of the depth. A differential equation for v
y can be obtained by writing dv dv dv dy = =v  dt dy dt dy This gives us the initial value problem √

dv 65 − 05v = dy 167v

10/3



v
0 = 0

(5.1)

This equation is also separable, but we cannot perform the integrations needed to find v
y explicitly. We have reached a position that is common in modeling a real-world phenomenon. The model (5.1) does not admit a closed form solution. At this point we will opt for a numerical approach to obtain approximate values for the velocity. But life is not this easy! If we attempt Euler’s method on the problem with equation (5.1), we cannot even get started because the initial condition is v
0 = 0, and v occurs in the denominator. There is a way around this difficulty. Reverse perspective and look for depth as a function of velocity, y
v. We will then calculate y
40, the depth when the velocity reaches 40 feet per second. Since the velocity is an increasing function of the depth, if y
40 > 300 feet, we will know that the barrel could not have achieved a velocity of 40 feet per second when it reached the bottom. If y
40 < 300, then we will know that when the drum hits bottom it was moving at more than 40 feet per second, hence is likely to rupture. Since dy/dv = 1/
dv/dy, the initial value problem for y
v is dy 167v √ =  dv 65 − 05v 10/3

y
0 = 0

√ Write 10/3 ≈ 1054 and apply Euler’s method with h = 1 and n = 40. We get y
40 ≈ 2682 feet. With h = 05 and n = 80 we get y
40 ≈ 2723 feet. Further reductions in step size will provide better accuracy. With h = 01 and n = 400 we get y
40 ≈ 2755 feet. Based on these numbers it would appear that the drum will exceed 40 feet per second when it has fallen 300 feet, hence has a good chance of leaking dangerous material. A more detailed analysis, using an error bound that we have not discussed, leads to the conclusion that the drum achieves a velocity of 40 feet per second somewhere between 272 and 279 feet, giving us confidence that it has reached this velocity by the time it lands on the ocean floor. This led to the conclusion that the plan for storing radioactive waste materials in drums on the ocean floor is too dangerous to be feasible.

SECTION 5.1

PROBLEMS

In each of Problems 1 through 6, generate approximate numerical solutions using h = 02 and twenty iterations,

then h = 01 and forty iterations, and finally h = 005 and eighty iterations. Graph the approximate solutions on the

CHAPTER 5

190

Numerical Approximation of Solutions y
= 1/x y
1 = 0. Use h = 001. Will this approximation be less than or greater than the actual value?

same set of axes. Also obtain error bounds for each case. In each of Problems 1 through 5, obtain the exact solution and graph it with the approximate solutions.

9. In the analysis of the radioactive waste disposal problem, how does the constant of proportionality for the drag on the drum affect the conclusion? √ Carry out the numerical analysis if the drag is √03v 10/3 , and again for the case that the drag is 08v 10/3 .

1. y
= y sin
x y
0 = 1


2. y = x + y y
1 = −3 3. y
= 3xy y
0 = 5 4. y
= 2 − x y
0 = 1 5. y
= y − cos
x y
1 = −2


6. y = x − y  y
0 = 4 2

7. Approximate e as follows. Use Euler’s method with h = 001 to approximate y
1, where y
x is the solution of y
= y y
0 = 1. Sketch a graph of the solution before applying Euler’s method and determine whether the approximate value obtained is less than or greater than the actual value. 8. Approximate ln
2 by using Euler’s method to approximate y
2, where y
x is the solution of

5.2

√ 10. Try exponents other than 10/3 for the velocity in the disposal problem to gauge the effect of this number on the conclusion. In particular, perform the analysis √ if the drag equals 05v (1 is slightly less than 10/3) and again for√a drag effect of 05v4/3 (4/3 is slightly greater than 10/3). 11. Suppose the drums are dropped over a part of the ocean having a depth of 340 feet. Will the drums be likely to rupture on impact with the ocean floor?

One-Step Methods Euler’s method is a one-step method because the approximation at xk+1 depends only on the approximation at xk , one step back. We will consider some other one-step methods for the initial value problem y
= f
x y

y
x0  = y0 

As usual, let the step size be h, and denote xk = x0 + kh for k = 0 1 2     n.

5.2.1

The Second-Order Taylor Method

By Taylor’s theorem with remainder (under certain conditions on f and h) we can write 1 2

h y
xk  2! 1 1 + · · · + hm y
m
xk  + hm+1 y
m+1
k  m!
m + 1!

y
xk+1  = y
xk  + hy

xk  +

for some k in xk  xk+1 . If y
m+1
x is bounded, then the last term in this sum can be made as small as we like by choosing h small enough. We therefore form the approximation yk+1 ≈ y
xk  + hy

xk  +

1 2

1 h y
xk  + · · · + hm y
m
xk  2! m!

If m = 1, this is Euler’s method, since y

xk  = f
xk  yk . Now let m = 2. Then yk+1 ≈ y
xk  + hy

xk  +

1 2

h y
xk  2!

(5.2)

5.2 One-Step Methods

191

We know that y
x = f
x y
x. This suggests that in the approximation (5.2) we consider f
xk  yk  as an approximation of y

xk  if yk is an approximation of y
xk . Thus consider y

xk  ≈ f
xk  yk  This leaves the term y


xk  in the approximation (5.2) to treat. First differentiate the expression y

x = f
x y
x to get y


x =

f

f
x y +
x yy

x

x

y

This suggests we consider y


xk  ≈

f

f
xk  yk  +
xk  yk f
xk  yk 

x

y

Insert these approximations of y

xk  and y


xk  into the approximation (5.2) to get  

f 1

f
xk  yk  +
xk  yk f
xk  yk   yk+1 ≈ yk + hf
xk  yk  + h2 2

x

y This is a one-step method, because yk+1 is obtained from information at xk , one step back from xk+1 

DEFINITION 5.2

Second-Order Taylor Method

The second-order Taylor method consists of approximating y
xk+1  by the expression   1 2 f

f yk+1 ≈ yk + hf
xk  yk  + h
x  y  +
xk  yk f
xk  yk   2

x k k

y

This expression can be simplified by adopting the notation fk = f
xk  yk 

f = fx 

x

f = fy 

y

and

f
x  y  =
fx k = fxk 

x k k

f
x  y  =
fy k = fyk 

y k k

Now the formula is  1  yk+1 ≈ yk + hfk + h2 fxk + fk fyk  2

192

CHAPTER 5

Numerical Approximation of Solutions

EXAMPLE 5.3

Consider y
= y2 cos
x

y
0 = 1/5

With f
x y = y2 cos
x, we have fx = −y2 sin
x and fy = 2y cos
x. Form 1 yk+1 ≈ yk + hyk2 cos
xk  + h2 yk2 cos2
xk  − h2 yk2 sin
xk  2 With h = 02 and twenty iterations (n = 20) we get the approximate values given in Table 5.4, for points xk = 0 + 02k for k = 0 1     20.

TA B L E 5 . 4

Approximation Values of the Solution of y
= y2 cos
x y
0 = 1/5

x

yapp
x

x

yapp
x

0.0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 1.6 1.8 2.0

0.2 0.20832 0.2170013470 0.2256558280 0.2337991830 0.2408797598 0.2463364693 0.2496815188 0.2505900093 0.2489684556 0.2449763987

2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

02389919589 02315347821 02231744449 02144516213 02058272673 01976613648 01902141527 01836603456 01781084317 01736197077

y 0.24

0.22

0.20

0.18 x 0

1

2

3

4

FIGURE 5.9 Exact and second-order Taylor approximate solutions of y
= y2 cos
x y
0 = 15 .

This problem can be solved exactly, and we obtain y
x = 1/
5 − sin
x. Figure 5.9 shows a graph of this solution, together with a smooth curve drawn through the approximated function values. The student should redo the approximation, using h = 01 and n = 40 for comparison. The Euler approximations for this example, with h = 02, are yk+1 = yk +
02yk2 cos
xk 

5.2 One-Step Methods

193

It is instructive to compute these approximations for n = 20, and compare the accuracy of the Euler method with that of the second-order Taylor method for this problem.

5.2.2

The Modified Euler Method

Near the end of the nineteenth century the German mathematician Karl Runge noticed a similarity between part of the formula for the second-order Taylor method and another Taylor polynomial approximation. Write the second-order Taylor formula as

1 yk+1 = yk + h fk + h
fx
xk  yk  + fk fy
xk  yk   (5.3) 2 Runge observed that the term in square brackets on the right side of this equation resembles the Taylor approximation f
xk + h yk + k ≈ fk + hfx
xk  yk  + hfy
xk  yk  In fact, the term in square brackets in equation (5.3) is exactly the right side of the last equation if we choose  =  = 1/2. This suggests the approximation   h hfk yk+1 ≈ yk + hf xk +  yk +  2 2

DEFINITION 5.3

Modified Euler Method

The modified Euler method consists of defining the approximation yk+1 by   h hfk yk+1 ≈ yk + hf xk +  yk +  2 2

The method is in the spirit of Euler’s method, except that f
x y is evaluated at
xk + h/2 yk + hfk /2 instead of at
xk  yk . Notice that
xk + h/2 is midway between xk and xk+1 .

EXAMPLE 5.4

Consider 1 y
− y = 2x2  x

y
1 = 4

Write the differential equation as y
=

y + 2x2 = f
x y x

Using the modified Euler method with h = 02 and n = 20 iterations, generate the approximate solution values given in Table 5.5.

194

CHAPTER 5

Numerical Approximation of Solutions

The exact solution of this initial value problem is y
x = x3 + 3x. The graph of this solution, together with a smooth curve drawn through the approximated values, are shown in Figure 5.10. The two curves coincide in the scale of the drawing. For example, y
5 = 140, while y20 , the approximated solution value at x20 = 1 + 20
02 = 5, is 1397. This small a difference does not show up on the graph. TA B L E 5 . 5

Approximate Values of the Solution of y
=
y/x + 2x2  y
1 = 4

x

yapp
x

x

yapp
x

1.0 1.2 1.4 1.6 1.8 2.0 2.2 2.4 2.6 2.8 3.0

4 5320363636 6927398601 8869292639 1119419064 1395020013 1718541062 2094789459 2525871247 3024691542 3587954731

3.2 3.4 3.6 3.8 4.0 4.2 4.4 4.6 4.8 5.0

4223164616 4935124526 5728637379 6608505841 7579532194 8646518560 9814266841 1108757877 1247125592 1397009975

y 140 120 100 80 60 40 20 0 1

2

3

4

5

x

FIGURE 5.10 Exact and modified Euler approximation of the solution of y
−
1/xy = 2x2  y
1 = 4.

We leave it for the student to do this example using the other approximation schemes for comparison. For this initial value problem, the other methods with h = 02 give:   1 2 Euler: yk+1 = yk + 02 y + 2xk  xk k  ⎤ ⎡ yk + 01 xyk + 2xk2 k + 2
xk + 012 ⎦  Modified Euler: yk+1 = yk + 02 ⎣ xk + 01 Second-order Taylor:

  yk yk yk 2 + 2xk + 002 − + 4xk + + 2xk  yk+1 = yk + 02 xk xk xk

5.2 One-Step Methods

5.2.3

195

Runge-Kutta Methods

An entire class of one-step methods is generated by replacing the right side in the modified Euler method with the general form afk + bf
xk + h yk + hfk  The idea is to choose the constants a, b,  and  to obtain an approximation with as favorable an error bound as possible. The fourth order Runge-Kutta method (known as RK4) has proved both computationally efficient and accurate, and is obtained by a clever choice of these constants in approximating slopes at various points. Without derivation, we will state the method.

DEFINITION 5.4

RK4

The RK4 method of approximation is to define yk+1 in terms of yk by yk+1 = yk +

h

W + 2Wk2 + 2Wk3 + Wk4   6 k1

where Wk1 = fk  Wk2 = f
xk + h/2 yk + hWk1 /2 Wk3 = f
xk + h/2 yk + hWk2 /2 Wk4 = f
xk + h yk + hWk3 

EXAMPLE 5.5

Consider y
=

1 cos
x + y y
0 = 1 y

Choose h = 01 and n = 40 to obtain approximate solution values at 1, 11 ,    , 39, 4. Compute Wk1 =

1 cos
xk + yk  yk

Wk2 =

1 cos
xk + yk + 005 + 005Wk1  yk + 005Wk1

Wk3 =

1 cos
xk + yk + 005 + 005Wk2  yk + 005Wk2

Wk4 =

1 cos
xk + yk + 01 + 01Wk3  yk + 01Wk3

and 01
Wk1 + 2Wk2 + 2Wk3 + Wk4  6 For comparison, we also compute approximations using modified Euler:   01 005 cos
xk + yk  yk+1 = yk +  cos xk + yk + 005 + yk + 005

1/yk  cos
xk + yk  yk yk+1 = yk +

196

CHAPTER 5

Numerical Approximation of Solutions

TA B L E 5 . 6

Approximate Values of the Solution of y
=
1/y cos
x + y y
0 = 1

x

RK4 Application

Modified Euler Application

x

RK4 Application

Modified Euler Application

0.0 0.1 0.2 0.3 0.4 0.5 0.6 0.7 0.8 0.9 1.0 2.2 2.3 2.4 2.5 2.6 2.7 2.8 2.9 3.0 3.1 3.2

1 1046496334 1079334533 1100247716 1110582298 1111396688 1103521442 1087597013 1064096638 1033337679 09954821879 −1711647122 −176180864 −1809127759 −1853613925 −1895247358 −1933982047 −1969747992 −2002452881 −2031983404 −2058206363 −2080969742

1 1046149156 1078757139 1099515885 1109747010 1110493110 1102574266 1086623840 1063110255 1032347712 09944964300 −1168193825 −1209281807 −1244631007 −1274170205 −1297734699 −1315071925 −1325843341 −1329623556 −1325897430 −1314055719 −1293389806

1.1 1.2 1.3 1.4 1.5 1.6 1.7 1.8 1.9 2.0 2.1 3.3 3.4 3.5 3.6 3.7 3.8 3.9 4.0

09505251773 08982661137 08382537806 07696832419 06911954958 06004493273 04930456632 03589118368 0156746754 −1602586399 −1658600340 −2100103902 −2115423092 −2126727488 −2133806037 −2136440412 −2134410465 −2127506124 −2115514719

09495508719 08973115249 0837306783 07688117656 06904155222 05998524927 04928907147 03602630404 01719452522 −1068545688 −1121333887 −1263086041 −1222220133 −1169751386 −1104515644 −1025210458 −0930599244 −08181811612 −06862852770

y 1.0 0.5 x 0

1

2

3

4

–0.5 –1.0

Euler modified

–1.5 –2.0

RK4

FIGURE 5.11 Runge-Kutta and modified Euler approximations of the solution of y
=
1/y cos
x + y y
0 = 1.

Table 5.6 shows the computed values, and Figure 5.11 shows graphs drawn through the approximated points by both methods. The two graphs are in good agreement as x nears 2, but then they diverge from each other. This divergence can be seen in the table. In general, RK4 is more accurate than modified Euler, particularly as the distance increases from the point where the initial data is specified. It can be shown that the Taylor and modified Euler methods are of order h2 , while RK4 is of order h4 . Since usually 0 < h < 1, h4 < h2 < h , so accuracy improves much faster by choosing smaller h with RK4 than with the other methods. There are higher order Runge-Kutta

5.3 Multistep Methods

197

methods which are of order hp for larger p. Such methods offer greater accuracy, but usually at a cost of more computing time. All of the methods of this section are one-step methods, which have the general form yk+1 = yk + 
xk  yk . In the next section, we will discuss multistep methods.

SECTION 5.2

PROBLEMS

In each of Problems 1 through 6, use the modified Euler, Taylor and RK4 methods to approximate the solution, first using h = 02 with twenty iterations, then h = 01 with forty iterations, and finally h = 005 with eighty iterations. Graph the approximate solutions for each method, and for a given h, on the same set of axes. 1. y
= sin
x + y y
0 = 2 2. y
= y − x2  y
1 = −4. Also solve this problem exactly and include a graph of the exact solution with graphs of the approximate solutions. 3. y
= cos
y + e−x  y
0 = 1 4. y
= y3 − 2xy y
3 = 2 5. y
= −y + e−x  y
0 = 4. Also solve this problem exactly and include a graph of the exact solution with graphs of the approximate solutions. 6. y
= sec
1/y − xy2  y
/4 = 1

5.3

7. Do Problem 3 of Section 5.1 with RK4 instead of Euler. Which method yields the better result? 8. Do Problem 5 of Section 5.1 with RK4 instead of Euler. Which method yields the better result? 9. Derive the improved Euler method, also known as the Heun method, as follows. Begin with Euler’s method and replace fk with
fk + fk+1 /2. Next replace yk+1 in fk+1 with yk + hfk . The result should be the approximation scheme yk+1 = yk +
h/2
fk + f
xk+1  yk + hfk . In each of Problems 10 through 12, use Euler, modified Euler and improved Euler to approximate the solution. Use h = 02 with n = 20, then h = 01 with n = 40 and then h = 005 with n = 80. Graph the approximate solutions, for each h, on the same set of axes. Whenever the solution can be found in exact form, graph this solution with the approximate solutions for comparison. 10. y
= 1 − y y
0 = 2 11. y
= −y/x + x y
1 = 1 12. y
= y − ex  y
−1 = 4

Multistep Methods We continue with the problem y
= f
x y y
x0  = y0 . The solution is y
x, and we presumably do not have a "good" expression for this function. We want to obtain approximate values yk to y
xk , where xk = x0 + kh for k = 0 1     n. The basis for some multistep methods is the informal belief that, if pk
x is a polynomial that x x approximates f
x y
x on xk  xk+1 , then xkk+1 pk
xdx should approximate xkk+1 f
x y
xdx. Now write
xk+1
xk+1
xk+1 y
xk+1  − y
xk  = y

xdx = f
x y
xdx ≈ pk
xdx xk

xk

Therefore y
xk+1  ≈ y
xk  +

xk




xk+1

xk

pk
xdx

(5.4)

So far this is deliberately vague. Nevertheless, this proposed approximation (5.4) contains the germ of an idea which we will now pursue.

198

CHAPTER 5

Numerical Approximation of Solutions

First we must decide how to choose the polynomials pk
x. Suppose we have somehow arrived at satisfactory approximations yk , yk−1      yk−r to the solution at xk , xk−1      xk−r , respectively. Then fk = f
xk  yk  ≈ f
xk  y
xk  fk−1 = f
xk−1  yk−1  ≈ f
xk−1  y
xk−1      fk−r = f
xk−r  yk−r  ≈ f
xk−r  y
xk−r  Keep in mind here that y
xk  is the exact solution evaluated at xk (this is unknown to us), and yk is an approximation of this solution value, obtained by some scheme. Now choose pk
x to be the polynomial of degree r passing through the points
xk  fk 
xk−1  fk−1     
xk−r  fk−r  These r +1 points will uniquely determine the r-degree polynomial pk
x. When this polynomial is inserted into the approximation scheme (5.4), we obtain a multistep approximation method in which the approximation yk+1 of y
xk+1  is defined by yk+1 = yk +




xk+1 xk

pk
xdx

(5.5)

We obtain different methods for different choices of r. Consider some cases of interest.

5.3.1

Case 1 r = 0

Now pk
x is a zero degree polynomial, or constant. Specifically, pk
x = fk for xk ≤ x ≤ xk+1 . The approximation scheme defined by equation (5.5) becomes yk+1 = yk +




xk+1 xk

fk dx = yk + fk xk+1 − xk  = yk + hfk

for k = 0 1 2    , n − 1. This is Euler’s method, a one-step method.

5.3.2

Case 2 r = 1

Now pk
x is a first degree polynomial, whose graph is the straight line through
xk  fk  and
xk−1  fk−1 . Therefore 1 1 pk
x = −
x − xk fk−1 +
x − xk−1 fk  h h Upon inserting this into the scheme (5.5) we get yk+1 = yk +




xk+1 xk



1 1 −
x − xk fk−1 +
x − xk−1 fk dx h h

5.3 Multistep Methods

199

A routine integration which we omit yields h yk+1 = yk + 3fk − fk−1  2 for k = 1 2     n − 1. This is a two-step method because computation of yk+1 requires prior computation of information at the two points xk and xk−1 . For larger r, the idea is the same, but the details are move involved because pk
x is of degree r and the integral in equation (5.5) is more involved, though still elementary. Here are the final results for two more cases.

5.3.3

Case 3 r = 3 yk+1 = yk +

h

23fk − 16fk−1 + 5fk−2  12

for k = 2 3     n − 1. This is a three-step method, requiring information at three points to compute the approximation yk+1 .

5.3.4

Case 4 r = 4 yk+1 = yk +

h

55fk − 59fk−1 + 37fk−2 − 9fk−3  24

(5.6)

for k = 3 4     n − 1. This is a four-step method. We might expect multistep methods to improve in accuracy as the number of steps increases, since more information is packed into the computation of the approximation at the next point. This is in general true, and an r-step method using an interpolating polynomial of degree r on each subinterval has error of order O
hr . The cost in improved accuracy is that the polynomials must be computed on each interval and more data is put into computation of each successive yk+1 . The schemes just given for r = 1 2, and 3 are called Adams-Bashforth multistep methods. One drawback to a multistep method is that some other method must be used to initiate it. For example, equation (5.6) involves fk−3 , and so is only valid for k = 3 4     n − 1. Some other scheme must be used to start it by feeding in y1 , y2 and y3 (and, of course, y0 is given as information). Often RK4 is used as an initiator in computing these first values. Another class of multistep methods, called Adams-Moulton methods, is obtained by using different data points to determine the interpolating polynomial pk
x to use in equation (5.4). For r = 2, pk
x is now chosen as the unique second degree polynomial passing through
xk−1  fk−1 ,
xk  fk  and
xk+1  fk+1 . This leads to the approximating scheme yk+1 = yk +

h

9f + 19fk − 5fk−1 + fk−2  24 k+1

(5.7)

This Adams-Moulton method is a four step method, and has error of order O
h4 . There is a significant difference between the Adams-Bashforth method (5.6) and the AdamsMoulton method (5.7). The former determines yk+1 in terms of three previously computed quantities, and is said to be explicit. The latter contains yk+1 on both sides of the equation (5.7), because yk+1 = f
xk+1  yk+1 , and therefore defines yk+1 implicitly by an equation containing yk+1 on both sides. Equation (5.7) therefore only provides an equation containing yk+1 , from which yk+1 must then be extracted, a perhaps nontrivial task.

CHAPTER 5

200

Numerical Approximation of Solutions

SECTION 5.3

PROBLEMS

In each of Problems 1 through 5, use the Taylor, modified Euler, and RK4 methods to approximate solution values. First use h = 02 with 20 iterations, then h = 01 with 40 iterations, then h = 005 with 80 iterations.


1. y = 4y − x y
3 = 0 2

2. y
= x sin
y − x2  y
1 = −3 3. y
= x2 + 4y y
0 = −2 4. y
= 1 − cos
x − y + x2  y
3 = 6 5. y
= 4x3 − xy + cos
y y
0 = 4 In each of Problems 6, 7 and 8, use the Adams-BashforthMoulton scheme, first with h = 02 and twenty iterations, then with h = 01 and forty iterations. 6. y
= y − x3  y
−2 = −4 7. y
= 2xy − y3  y
0 = 2 8. y
= ln
x + x2 y y
2 = 1

9. Carry out the details for deriving the two-step scheme stated for the case r = 2. 10. Carry out the details for deriving the three-step scheme stated for the case r = 3. 11. Every one-step and multistep method we have considered is a special case of the general expression yk+1 =

n 

j yk+1−j

j=1

+ h
xk+1−m      xk  xk+1  yk+1−m      yk  yk+1  By making appropriate choices of m, the 
j s and , show how this formula gives Euler’s method, the modified Euler method, the Taylor method, RK4, and the Adams-Bashforth method.

PA RT

2

CHAPTER 6 Vectors and Vector Spaces

CHAPTER 7

Vectors and Linear Algebra

Matrices and Systems of Linear Equations

CHAPTER 8 Determinants

CHAPTER 9 Eigenvalues, Diagonalization, and Special Matrices

Some quantities are completely determined by their magnitude, or “size.” This is true of temperature and mass, which are numbers referred to some scale or measurement system. Such quantities are called scalars. Length, volume, and distance are other scalars. By contrast, a vector carries with it a sense of both magnitude and direction. The effect of a push against an object will depend not only on the magnitude or strength of the push, but also on the direction in which it is exerted. This part is concerned with the notation and algebra of vectors and objects called matrices. This algebra will be used to solve systems of linear algebraic equations and systems of linear differential equations. It will also give us the machinery needed for the quantitative study of systems of differential equations (Part 3), in which we attempt to determine the behavior and properties of solutions when we cannot write these solutions explicitly or in closed form. In Part 4, vector algebra will be used to develop vector calculus, which extends derivatives and integrals to higher dimensions, with applications to models of physical systems, partial differential equations, and the analysis of complex-valued functions.

201

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CHAPTER

6

ALGEBRA AND GEOMETRY OF VECTORS THE DOT PRODUCT THE CROSS PRODUCT THE VECTOR SPACE R n LINEAR INDEPENDENCE SPANNING SETS AND DIMENSION IN R n ABSTRACT VECTOR SPACES

Vectors and Vector Spaces

6.1

The Algebra and Geometry of Vectors When dealing with vectors, a real number is often called a scalar. The temperature of an object and the grade of a motor oil are scalars. We want to define the concept of a vector in such a way that the package contains information about both direction and magnitude. One way to do this is to define a vector (in 3-dimensional space) as an ordered triple of real numbers.

DEFINITION 6.1

Vector

A vector is an ordered triple
a b c, in which a, b, and c are real numbers.

We represent the vector
a b c as an arrow from the origin
0 0 0 to the point
a b c in 3-space, as in Figure 6.1. In this way, the direction indicated by the arrow, as viewed from the origin, gives the direction of the vector. The length of the arrow is the magnitude (or norm) of the vector—a longer arrow represents a vector of greater strength. Since the distance from √ the origin to the point
a b c is a2 + b2 + c2 , we will define this number to be the magnitude of the vector
a b c.

DEFINITION 6.2

Norm of a Vector

The norm, or magnitude, of a vector
a b c is the number 
a b c defined by  
a b c = a2 + b2 + c2 203

204

CHAPTER 6

Vectors and Vector Spaces z

z

z

(a, b, c)

18

(a, b, c)

(
1, 4, 1)

y x

y

y x

FIGURE 6.1 The vector
a b c is represented by the arrow from
0 0 0 to the point
a b c.

x

FIGURE 6.2


−1 4 1 =

FIGURE 6.3 Parallel representations of the same vector.

√ 18.

√ √ For example, the norm of
−1 4 1 is 
−1 4 1 = 1 + 16 + 1 = 18. This is the length of the arrow from the origin to the point
−1 4 1 (Figure 6.2). The only vector that is not represented by an arrow from the origin is the zero vector
0 0 0, which has zero magnitude and no direction. It is, however, useful to have a zero vector, because various forces in a physical process may cancel each other, resulting in a zero force or vector. The number a is the first component of
a b c, b is the second component, and c the third component. Two vectors are equal if and only if each of their respective components is equal:
a b c =
u v w if and only if a = u

b = v

c = w

We will usually denote scalars (real numbers) by letters in regular type face (a, b, c, A, B,    ), and vectors by letters in boldface (a, b, c, A, B,    ). The zero vector is denoted O. Although there is a difference between a vector (ordered triple) and an arrow (visual representation of a vector), we often speak of vectors and arrows interchangeably. This is useful in giving geometric interpretations to vector operations. However, any two arrows having the same length and same direction are said to represent the same vector. In Figure 6.3, all the arrows represent the same vector. We will now develop algebraic operations with vectors and relate them to the norm.

DEFINITION 6.3

Product of a Scalar and Vector

The product of a real number  with a vector F =
a b c is denoted F, and is defined by F =
a b c

Thus a vector is multiplied by a scalar by multiplying each component by the scalar. For example, 3
2 −5 1 =
6 −15 3

and

− 5
−4 2 10 =
20 −10 −50

The following relationship between norm and the product of a scalar with a vector leads to a simple geometric interpretation of this operation.

6.1 The Algebra and Geometry of Vectors

205

THEOREM 6.1

Let F be a vector and  a scalar. Then 1. F =  F. 2. F = 0 if and only if F = O. Proof

If F =
a b c, then F =
a b c, so  F = 2 a2 + 2 b2 + 2 c2  =  a2 + b2 + c2 =  F 

This proves conclusion (1). For (2), first recall that O =
0 0 0, so  O = 02 + 02 + 02 = 0 Conversely, if F = 0, then a2 + b2 + c2 = 0, hence a = b = c = 0 and F = O Consider this product of a scalar with a vector from a geometric point of view. By (1) of the theorem, the length of F is  times the length of F. Multiplying by  lengthens the arrow representing F if  > 1, and shrinks it to a shorter arrow if 0 <  < 1. Of course, if  = 0 then F is the zero vector, with zero length. But the algebraic sign of  has an effect as well. If  is positive, then F has the same direction as F, while if  is negative, F has the opposite direction.

EXAMPLE 6.1

Let F =
2 4 1, as shown in Figure 6.4. 3F =
6 12 3 is along the same direction as F, but is represented as an arrow three times longer. But −3F =
−6 −12 −3, while being three times longer than F, is in the direction opposite that of F through the origin. And 21 F =
1 2 1/2 is in the same direction as F, but half as long.

z (6, 12, 3) (1, 2, 12 ) (2, 4, 1) y (
6,
12,
3)

x

FIGURE 6.4 Scalar multiples of a vector.

In particular, the scalar product of −1 with F =
a b c is the vector
−a −b −c, having the same length as F, but the opposite direction. This vector is called “minus F,” or the negative of F, and is denoted −F. Consistent with the interpretation of multiplication of a vector by a scalar, we define two vectors F and G to be parallel if each is a nonzero scalar multiple of the other. Of course if F = G and  = 0, then G =
1/F. Parallel vectors may differ in length, and even be in opposite directions, but the straight lines through arrows representing these vectors are parallel lines in 3-space.

206

CHAPTER 6

Vectors and Vector Spaces

The algebraic sum of two vectors is defined as follows.

DEFINITION 6.4

Vector Sum

The sum of F =
a1  b1  c1  and G =
a2  b2  c2  is the vector F + G =
a1 + a2  b1 + b2  c1 + c2 

That is, we add vectors by adding respective components. For example,
−4  2 +
16 1 −5 =
12  + 1 −3 If F =
a1  b1  c1  and G =
a2  b2  c2 , then the sum of F with −G is
a1 − a2  b1 − b2  c1 − c2 . It is natural to denote this vector as F − G, and refer to it as “F minus G.” For example,
−4  2 minus
16 1 −5 is
−4  2 −
16 1 −5 =
−20  − 1 7 We therefore subtract two vectors by subtracting their respective components. Vector addition and multiplication of a vector by a scalar have the following computational properties. THEOREM 6.2

Algebra of Vectors

Let F, G, and H be vectors and let  and  be scalars. Then 1. 2. 3. 4. 5. 6.

F + G = G + F.
F + G + H = F +
G + H. F + O = F. 
F + G = F + G.
F = 
F
 + F = F + F

Conclusion (1) is the commutative law for vector addition, and (2) is the associative law. Conclusion (3) states that the zero vector behaves with vectors like the number zero does with real numbers, as far as addition is concerned. The theorem is proved by routine calculations, using the properties of real-number arithmetic. For example, to prove (4), write F =
a1  b1  c1  and G =
a2  b2  c2 . Then 
F + G = 
a1 + a2  b1 + b2  c1 + c2  =

a1 + a2  
b1 + b2  
c1 + c2  =
a1 + a2  b1 + b2  c1 + c2  =
a1  b1  c1  +
a2  b2  c2  = 
a1  b1  c1  + 
a2  b2  c2  = F + G Vector addition has a simple geometric interpretation. If F and G are represented as arrows from the same point P, as in Figure 6.5, then F + G is represented as the arrow from P to the opposite vertex of the parallelogram having F and G as two incident sides. This is called the parallelogram law for vector addition.

6.1 The Algebra and Geometry of Vectors

207

z (0, 0, 1)

F

G

F G

G

G G

P

(0, 1, 0)

F

F

k

G

y

j

i (1, 0, 0)

F

F

x

FIGURE 6.5

FIGURE 6.6

Parallelogram law for vector addition.

Another way of visualizing the parallelogram law.

FIGURE 6.7 Unit vectors along the axes.

The parallelogram law suggests a strategy for visualizing addition that is sometimes useful. Since two arrows having the same direction and length represent the same vector, we could apply the parallelogram law to form F + G as in Figure 6.6, in which the arrow representing G is drawn from the tip of F, rather from a common initial point with F. We often do this in visualizing computations with vectors. Any vector can be written as a sum of scalar multiples of “standard” vectors as follows. Define i =
1 0 0

j =
0 1 0

k =
0 0 1

These are unit vectors (length 1) aligned along the three coordinate axes in the positive direction (Figure 6.7). In terms of these vectors, F =
a b c = a
1 0 0 + b
0 1 0 + c
0 0 1 = ai + bj + ck This is called the standard representation of F. When a component of F is zero, we usually just omit it in the standard representation. For example,
−3 0 1 = −3i + k Figure 6.8 shows two points P1
a1  b1  c1  and P2
a2  b2  c2 . It will be useful to know the vector represented by the arrow from P1 to P2 . Let H be this vector. Denote G = a1 i + b1 j + c1 k

and F = a2 i + b2 j + c2 k

By the parallelogram law (Figure 6.9), G + H = F Hence H = F − G =
a2 − a1 i +
b2 − b1 j +
c2 − c1 k For example, the vector represented by the arrow from
−2 4 1 to
14 5 −7 is 16i + j − 8k. The vector from
14 5 −7 to
−2 4 1 is the negative of this, or −16i − j + 8k. Vector notation and algebra are often useful in solving problems in geometry. This is not our goal here, but the reasoning involved is often useful in solving problems in the sciences and engineering. We will give three examples to demonstrate the efficiency of thinking in terms of vectors.

208

CHAPTER 6

Vectors and Vector Spaces (a 1, b1, c1 )

z

z

H  (a 2
a 1, b 2
b1, c2
c1 )

G

(a 1, b1, c1)

(a 2, b2, c2)

(a 2, b2, c2)

F  (a 2, b2, c2)

y

y x

x

FIGURE 6.8

FIGURE 6.9

The arrow from
a1  b1  c1  to
a2  b2  c2  is
a2 − a1 i +
b2 − b1 j +
c2 − c1 k.

EXAMPLE 6.2

Suppose we want the equation of the line L through the points
1 −2 4 and
6 2 −3. This problem is more subtle in 3-space than in the plane because in three dimensions there is no point-slope formula. Reason as follows. Let
x y z be any point on L. Then (Figure 6.10), the vector represented by the arrow from
1 −2 4 to
x y z must be parallel to the vector from
1 −2 4 to
6 2 −3, because arrows representing these vectors are both along L. This means that
x − 1i +
y + 2j +
z − 4k is parallel to 5i + 4j − 7k. Then, for some scalar t,
x − 1i +
y + 2j +
z − 4k = t 5i + 4j − 7k But then the respective components of these vectors must be equal: x − 1 = 5t

y + 2 = 4t

z − 4 = −7t

x = 1 + 5t

y = −2 + 4t

Then z = 4 − 7t

(6.1)

A point is on L if and only if its coordinates are
1 + 5t −2 + 4t 4 − 7t for some real number t (Figure 6.11). Equations (6.1) are parametric equations of the line, with t, which can be assigned any real value, as parameter. When t = 0, we get
1 −2 4, and when t = 1, we get
6 2 −3. We can also write the equation of this line in what is called normal form. By eliminating t, this form is x−1 y+2 z−4 = =  5 4 −7 We may also envision the line as swept out by the arrow pivoted at the origin and extending to the point
1 + 5t −2 + 4t 4 − 7t as t varies over the real numbers. Some care must be taken in writing the normal form of a straight line. For example, the line through
2 −1 6 and
−4 −1 2 has parametric equations x = 2 − 6t

y = −1

z = 6 − 4t

If we eliminate t, we get x−2 z−6 =  −6 −4

y = −1

6.1 The Algebra and Geometry of Vectors

209

z

(x, y, z)

(x
1)i  (y  2)j  (z
4)k (1,
2, 4) L y

5i  4j
7k

z

(1  5t,
2  4t, 4
7t)

L y

x ( 6, 2,
3)

x

FIGURE 6.10

FIGURE 6.11

Every point on the line has second coordinate −1, and this is independent of t. This information must not be omitted from the equations of the line. If we omit y = −1 and write just x−2 z−6 =  −6 −4 then we have the equation of a plane, not a line.

EXAMPLE 6.3

Suppose we want a vector F in the x y−plane, making an angle of /7 with the positive x−axis, and having magnitude 19. By “find a vector” we mean determine its components. Let F = ai + bj. From the right triangle in Figure 6.12, cos
/7 =

a 19

and

sin
/7 =

b  19

Then F = 19 cos
/7i + 19 sin
/7j

y 19 π 7

a

b x

FIGURE 6.12

EXAMPLE 6.4

We will prove that the line segments formed by connecting successive midpoints of the sides of a quadrilateral form a parallelogram. Again, our overall objective is not to prove theorems of geometry, but this argument is good practice in the use of vectors. Figure 6.13 illustrates what we want to show. Draw the quadrilateral again, with arrows (vectors) A, B, C, and D as sides. The vectors x, y, u, and v drawn with dashed lines connect

210

CHAPTER 6

Vectors and Vector Spaces P1 B y x A

u

v

P0 FIGURE 6.13

C

D

FIGURE 6.14

the midpoints of successive sides (Figure 6.14). We want to show that x and u are parallel and of the same length, and that y and v are also parallel and of the same length. From the parallelogram law for vector addition and the definitions of x and u, 1 1 x = A+ B 2 2 and 1 1 u = C + D 2 2 But also by the parallelogram law, A + B is the arrow from P0 to P1 , while C + D is the arrow from P1 to P0 . These arrows have the same length, and opposite directions. This means that A + B = −
C + D But then x = −u, so these vectors are parallel and of the same length (just opposite in direction). A similar argument shows that y and v are also parallel and of the same length, completing the proof.

SECTION 6.1

PROBLEMS

In each of Problems 1 through 5, compute F + G, F − G, F, G, 2F, and 3G. √ 1. F = 2i − 3j + 5k G = 2i + 6j − 5k 2. F = i − 3k G = 4j

9. F = i − 2j G = i − 3j 10. F = −i + 4j G = −2i − 2j In each of Problems 11 through 15, determine F and represent F and F as arrows.

3. F = 2i − 5j G = i + 5j − k √ 4. F = 2i + j − 6k G = 8i + 2k

11. F = i + j  = −1/2

5. F = i + j + k G = 2i − 2j + 2k

13. F = −3j  = −4

In each of Problems 6 through 10, calculate F + G and F − G by representing the vectors as arrows and using the parallelogram law.

14. F = 6i − 6j  = 1/2

6. F = i G = 6j 7. F = 2i − j G = i − j

In each of Problems 16 through 21, find the parametric equations of the straight line containing the given points. Also find the normal form of this line.

8. F = −3i + j G = 4j

16. (1,0,4),(2,1,1)

12. F = 6i − 2j  = 2

15. F = −3i + 2j  = 3

6.2 The Dot Product 17.
3 0 0
−3 1 0

25. 15 7/4

18.
2 1 1
2 1 −2

26. 25 3/2

19. (0,1,3), (0,0,1) 20.
1 0 −4
−2 −2 5 21.
2 −3 6
−1 6 4 In each of Problems 22 through 26, find a vector F in the x y-plane having the given length and making the angle (given in radians) with the positive x-axis. Represent the vector as an arrow in the plane. √ 22. 5 /4 23. 6 /3 24. 5 3/5

6.2

211

27. Let P1  P2      Pn be distinct points in 3-space, with n ≥ 3. Let Fi be the vector represented by the arrow from Pi to Pi+1 for i = 1 2     n − 1 and let Fn be the vector represented by the arrow from Pn to P1 . Prove that F1 + F2 + · · · + Fn = O. 28. Let F be any nonzero vector. Determine a scalar t such that tF = 1. 29. Use vectors to prove that the altitudes of any triangle intersect in a single point. (Recall that an altitude is a line from a vertex, perpendicular to the opposite side of the triangle.)

The Dot Product Throughout this section, let F = a1 i + b1 j + c1 k and G = a2 i + b2 j + c2 k

DEFINITION 6.5

Dot Product

The dot product of F and G is the number F · G defined by F · G = a1 a 2 + b 1 b 2 + c 1 c 2 

For example,

√ √
3i + 4j − k ·
−2i + 6j + 3k = −2 3 + 24 − 3

Sometimes the dot product is referred to as a scalar product, since the dot product of two vectors is a scalar (real number). This must not be confused with the product of a vector with a scalar. Here are some rules for operating with the dot product. THEOREM 6.3

Properties of the Dot Product

Let F, G, and H be vectors, and  a, scalar. Then 1. 2. 3. 4.

F · G = G · F.
F + G · H = F · H + G · H. 
F · G =
F · G = F ·
G. F · F = F2 .

5. F · F = 0 if and only if F = O. Conclusion (1) is the commutativity of the dot product (we can perform the operation in either order), and (2) is a distributive law. Conclusion (3) states that a constant factors through a

212

CHAPTER 6

Vectors and Vector Spaces

dot product. Conclusion (4) is very useful in some kinds of calculations, as we will see shortly. A proof of the theorem involves routine calculations, two of which we will illustrate. Proof

For (3), write 
F · G = 
a1 a2 + b1 b2 + c1 c2  =
a1 a2 +
b1 b2 +
c1 c2 =
F · G = a1
a2  + b1
b2  + c1
c2  = F ·
G

For (4), we have F · F =
a1 i + b1 j + c1 k ·
a1 i + b1 j + c1 k = a21 + b12 + c12 = F2  Using conclusion (4) of the theorem, we can derive a relationship we will use frequently. LEMMA 6.1

Let F and G be vectors, and let  and  be scalars. Then F + G2 = 2 F2 + 2F · G + 2 G  Proof

By using Theorem 6.3, we have F + G2 =
F + G ·
F + G = 2 F · F + F · G + G · F + 2 G · G = 2 F · F + 2F · G + 2 G · G = F2 + 2F · G + G2 

The dot product can be used to determine the angle between vectors. Represent F and G as arrows from a common point, as in Figure 6.15. Let  be the angle between F and G. The arrow from the tip of F to the tip of G represents G − F, and these three vectors form the sides of a triangle. Now recall the law of cosines, which states, for the triangle of Figure 6.16, that a2 + b2 − 2ab cos
 = c2 

(6.2)

Apply this to the vector triangle of Figure 6.15, with sides of length a = G, b = F, and c = G − F. By using Lemma 6.1 with  = −1 and  = 1, equation (6.2) becomes G2 + F2 − 2 G F cos
 = G − F2 = G2 + F2 − 2G · F Then F · G = F G cos
 Assuming that neither F nor G is the zero vector, then cos
 =

F·G  F G

(6.3)

This provides a simple way of computing the cosine of the angle between two arrows representing vectors. Since vectors can be drawn along straight lines, this also lets us calculate the angle between two intersecting lines.

6.2 The Dot Product

b c

θ

z (
1, 3, 1)

θ

G–F

G

213

y

θ

a

x (0, 2,
4)

F FIGURE 6.15

FIGURE 6.16

FIGURE 6.17

Law of cosines: a2 + b2 − 2ab cos
 = c2 .

EXAMPLE 6.5

Let F = −i + 3j + k and G = 2j − 4k. The cosine of the angle between these vectors (Figure 6.17) is cos
 = = Then


−i + 3j + k ·
2j − 4k −i + 3j + k 2j − 4k
−1
0 +
3
2 +
1
−4 2  =√ √ √ 2 2 2 2 2 1 +3 +1 2 +4 220 √  = arccos
2/ 220

√ which is that unique number in 0  whose cosine is 2/ 220.  is approximately 1436 radians.

EXAMPLE 6.6

Lines L1 and L2 are given, respectively, by the parametric equations x = 1 + 6t

y = 2 − 4t

z = −1 + 3t

and x = 4 − 3p

y = 2p

z = −5 + 4p

in which the parameters t and p take on all real values. We want the angle between these lines at their point of intersection, which is
1 2 −1 (on L1 for t = 0 and on L2 for p = 1). Of course, two intersecting lines have two angles between them (Figure 6.18). However, the sum of these angles is , so either angle determines the other. The strategy for solving this problem is to find a vector along each line, then find the angle between these vectors. For a vector F along L1 , choose any two points on this line, say
1 2 −1 and, with t = 1,
7 −2 2. The vector from the first to the second point is F =
7 − 1i +
−2 − 2j +
2 −
−1k = 6i − 4j + 3k

214

CHAPTER 6

Vectors and Vector Spaces z

y F

(7,
2, 2) L1

ϕ

θ

(1, 2,
1)

G

x

θ (4, 0,
5 ) L2

FIGURE 6.18

FIGURE 6.19

Two points on L2 are
1 2 −1 and, with p = 0,
4 0 −5. The vector G from the first to the second of these points is G =
4 − 1i +
0 − 2j +
−5 −
−1k = 3i − 2j − 4k These vectors are shown in Figure 6.19. The cosine of the angle between F and G is F·G 14 6
3 − 4
−2 + 3
−4 =√ =√  √ F G 36 + 16 + 9 9 + 4 + 16 1769 √ One angle between the lines is  = arccos
14/ 1769, approximately 123 radians. cos
 =

If we had √ used −G in place of G in this calculation, we would have gotten  = arccos
−14/ 1769, or about 191 radians. This is the supplement of the angle found in the example.

EXAMPLE 6.7

The points A
1 −2 1, B
0 1 6, and C
−3 4 −2 form the vertices of a triangle. Suppose we want the angle between the line AB and the line from A to the midpoint of BC. This line is a median of the triangle and is shown in Figure 6.20 Visualize the sides of the triangle as vectors, as in Figure 6.21. If P is the midpoint of BC, then H1 = H2 because both vectors have the same direction and length. From the coordinates of the vertices, calculate F = −i + 3j + 5k and G = −4i + 6j − 3k We want the angle between F and K, so we need K. By the parallelogram law, F + H1 = K and K + H2 = G Since H1 = H2 , these equations imply that K = F + H1 = F +
G − K Therefore, 1 5 9 K =
F + G = − i + j + k 2 2 2

6.2 The Dot Product A

A

θ

θ

B

215

G

F C

P

FIGURE 6.20

B

K H1

P

H2

C

FIGURE 6.21

Now the cosine of the angle we want is cos
 =

F·K 42 42 =√ √ =√  F K 35 110 3850

 is approximately 083 radians. The arrows representing two nonzero vectors F and G are perpendicular exactly when the cosine of the angle between them is zero, and by equation (6.3) this occurs when F · G = 0. This suggests we use this condition to define orthogonality (perpendicularity) of vectors. If we agree to the convention that the zero vector is orthogonal to every vector, then this dot product condition allows a general definition without requiring that the vectors be nonzero.

DEFINITION 6.6

Orthogonal Vectors

Vectors F and G are orthogonal if and only if F · G = 0.

EXAMPLE 6.8

Let F = −4i + j + 2k, G = 2i + 4k and H = 6i − j − 2k. Then F · G = 0, so F and G are orthogonal. But F · H and G · H are nonzero, so F and H are not orthogonal, and G and H are not orthogonal. Sometimes orthogonality of vectors is a useful device for dealing with lines and planes in three-dimensional space.

EXAMPLE 6.9

Two lines are given parametrically by L1  x = 2 − 4t

y = 6 + t

z = 3t

and L2  x = −2 + p

y = 7 + 2p

z = 3 − 4p

We want to know whether these lines are perpendicular. (It does not matter whether the lines intersect).

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The idea is to form a vector along each line and test these vectors for orthogonality. For a vector along L1 , choose two points on this line, say
2 6 0 when t = 0 and
−2 7 3 when t = 1. Then F = −4i + j + 3k is along L1 . Two points on L2 are
−2 7 3 when p = 0 and
−1 9 −1 when p = 1. Then G = i + 2j − 4k is along L2 . Since F · G = −14 = 0, these vectors, hence these lines, are not orthogonal.

EXAMPLE 6.10

Suppose we want the equation of a plane to the vector N = −2i + 4j + k.

containing the point
−6 1 1 and perpendicular

(
6, 1, 1) z

(x  6)i  (y
1) j  (z
1)k (x, y, z)
2i  4 j  k y

x FIGURE 6.22

A strategy to find such an equation is suggested by Figure 6.22. A point
x y z is on if and only if the vector from
−6 1 1 to
x y z is in and therefore is orthogonal to N. This means that

x + 6i +
y − 1j +
z − 1k · N = 0 Carrying out this dot product, we get the equation −2
x + 6 + 4
y − 1 +
z − 1 = 0 or −2x + 4y + z = 17 This is the equation of the plane. Of course the given point
−6 1 1 satisfies this equation. We will conclude this section with the important Cauchy-Schwarz inequality, which states that the dot product of two vectors cannot be greater in absolute value than the product of the lengths of the vectors. THEOREM 6.4

Cauchy-Schwarz Inequality

Let F and G be vectors. Then F · G ≤ F G  If either vector is the zero vector, then both sides of the proposed inequality are zero. Thus suppose neither vector is the zero vector. In this event,

Proof

cos
 =

F·G  F G

6.3 The Cross Product

217

where  is the angle between F and G. But then −1 ≤

F·G ≤ 1 F G

so − F G ≤ F · G ≤ F G  which is equivalent to the Cauchy-Schwarz inequality.

SECTION 6.2

PROBLEMS

In each of Problems 1 through 6, compute the dot product of the vectors and the cosine of the angle between them. Also determine if they are orthogonal and verify the Cauchy-Schwarz inequality for these vectors. 1. i 2i − 3j + k 2. 2i − 6j + k i − j

11.
0 −1 4 7i + 6j − 5k 12.
−2 1 −1 4i + 3j + k In each of Problems 13 through 16, find the cosine of the angle between AB and the line from A to the midpoint of BC. 13. A =
1 −2 6 B =
3 0 1 C =
4 2 −7

3. −4i − 2j + 3k 6i − 2j − k

14. A =
3 −2 −3 B =
−2 0 1 C =
1 1 7

4. 8i − 3j + 2k −8i − 3j + k

15. A =
1 −2 6 B =
0 4 −3 C =
−3 −2 7

5. i − 3k 2j + 6k

16. A =
0 5 −1 B =
1 −2 5 C =
7 0 −1

6. i + j + 2k i − j + 2k In each of Problems 7 through 12, find the equation of the plane containing the given point and having the given vector as normal vector. 7.
−1 1 2 3i − j + 4k

17. Suppose F · X = 0 for every vector X. What can be concluded about F? 18. Suppose F · i = F · j = F · k = 0. What can be concluded about F?

8.
−1 0 0 i − 2j

19. Suppose F = O. Prove that the unit vector u for which F · u is a maximum must be parallel to F.

9.
2 −3 4 8i − 6j + 4k

20. Prove that for any vector F, F =
F · ii +
F · jj +
F · kk

10.
−1 −1 −5 −3i + 2j

6.3

The Cross Product The dot product produces a scalar from two vectors. We will now define the cross product, which produces a vector from two vectors. For this section, let F = a1 i + b1 j + c1 k and G = a2 i + b2 j + c2 k.

DEFINITION 6.7

Cross Product

The cross product of F with G is the vector F × G defined by F × G =
b1 c2 − b2 c1 i +
a2 c1 − a1 c2 j +
a1 b2 − a2 b1 k

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This vector is read “F cross G.” For example,
i + 2j − 3k ×
−2i + j + 4k =
8 + 3i +
6 − 4j +
1 + 4k = 11i + 2j + 5k A cross product is often computed as a three by three “determinant,” with the unit vectors in the first row, components of F in the second row, and components of G in the third. If expanded by the first row, this determinant gives F × G. For example   i   1   −2

j 2 1

k −3 4

     2 =   1 

−3 4

    1 i−   −2

−3 4

    j+ 1   −2

2 1

  k 

= 11i + 2j + 5k = F × G The interchange of two rows in a determinant results in a change of sign. This means that interchanging F and G in the cross product results in a change of sign: F × G = −G × F This is also apparent from the definition. Unlike addition and multiplication of real numbers and the dot product operation, the cross product is not commutative, and the order in which it is performed makes a difference. This is true of many physical processes, for example, the order in which chemicals are combined may make a significant difference. Some of the rules we need to compute with cross products are given in the next theorem.

THEOREM 6.5

Properties of the Cross Product

Let F, G, and H be vectors and let  be a scalar. 1. F × G = −G × F. 2. F × G is orthogonal to both F and G. 3. F × G = F G sin
, in which  is the angle between F and G. 4. If F and G are not zero vectors, then F × G = O if and only F and G are parallel. 5. F ×
G + H = F × G + F × H. 6. 
F × G =
F × G = F ×
G.

Proofs of these statements are for the most part routine calculations. We will prove (2) and (3). Proof

For (2), compute F ·
F × G = a1
b1 c2 − b2 c1  + b1
a2 c1 − a1 c2  + c1
a1 b2 − a2 b1  = 0

Therefore F and F × G are orthogonal. A similar argument holds for G.

6.3 The Cross Product

219

For (3), compute F × G2 =
b1 c2 − b2 c1 2 +
a2 c1 − a1 c2 2 +
a1 b2 − a2 b1 2 =
a21 + b12 + c12 
a22 + b22 + c22  −
a1 a2 + b1 b2 + c1 c2 2 = F2 G2 −
F · G2 = F2 G2 − F2 G2 cos2
   = F2 G2 1 − cos2
 = F2 G2 sin2
 Because 0 ≤  ≤ , all of the factors whose squares appear in this equation are nonnegative, and upon taking square roots we obtain conclusion (3). If F and G are nonzero and not parallel, then arrows representing these vectors determine a plane in 3-dimensional space (Figure 6.23). F × G is orthogonal to this plane and oriented as in Figure 6.24. If a person’s right hand is placed so that the fingers curl from F to G, then the thumb points up along F × G. This is referred to as the right-hand rule. G × F = −F × G points in the opposite direction. As a simple example, i × j = k, and these three vectors define a standard right-handed coordinate system in 3-space.

G

G F

F

FIGURE 6.23

FIGURE 6.24

Plane determined by F and G.

Right-hand rule gives the direction of F × G.

The fact that F × G is orthogonal to both F and G is often useful. If they are not parallel, then vectors F and G determine a plane (Figure 6.23). This is consistent with the fact that three points, not on the same straight line, determine a plane. One point forms a base point for drawing the arrows representing F and G, and the other two points are the terminal points of these arrows. If we know a vector orthogonal to both F and G, then this vector is orthogonal to every vector in . Such a vector is said to be normal to . In Example 6.10 we showed how to find the equation of a plane, given a point in the plane and a normal vector. Now we can find the equation of a plane, given three points in it (not all on a line), because we can use the cross product to produce a normal vector.

EXAMPLE 6.11

Suppose we want the equation of the plane containing the points
1 2 1,
−1 1 3, and
−2 −2 −2. Begin by finding a vector normal to . We will do this by finding two vectors in and taking their cross product. The vectors from
1 2 1 to the other two given points are in (Figure 6.25). These vectors are F = −2i − j + 2k and G = −3i − 4j − 3k

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Vectors and Vector Spaces z

(
1, 1, 3) F

x

NF × G (1, 2, 1) y G (
2,
2,
2)

FIGURE 6.25

Form

  i j k  2 N = F × G =  −2 −1  −3 −4 −3

    = 11i − 12j + 5k  

This vector is normal to (orthogonal to every vector lying in ). Now proceed as in Example 6.10. If
x y z is any point in , then
x − 1i +
y − 2j +
z − 1k is in and so is orthogonal to N. Therefore,


x − 1i +
y − 2j +
z − 1k · N = 11
x − 1 − 12
y − 2 + 5
z − 1 = 0 This gives 11x − 12y + 5z = −8 This is the equation of the plane in the sense that a point
x y z is in the plane if and only its coordinates satisfy this equation. If we had specified three points lying on a line (collinear) in this example, then we would have found that F and G are parallel, hence F × G = O. When we calculated this cross product and got a nonzero vector, we knew that the points were not collinear. The cross product also has geometric interpretations as an area or volume.

THEOREM 6.6

Let F and G be represented by arrows lying along incident sides of a parallelogram (Figure 6.26). Then the area of this parallelogram is F × G. The area of a parallelogram is the product of the lengths of two incident sides and the sine of the angle between them. Draw vectors F and G along two incident sides. Then

Proof

G

θ FIGURE 6.26

= F × G.

F Area

6.3 The Cross Product

221

these sides have length F and G. If  is the angle between them, then the area of the parallelogram is F G sin
. But this is exactly F × G.

EXAMPLE 6.12

A parallelogram has two sides extending from
0 1 −2 to
1 2 2 and from
0 1 −2 to
1 4 1. We want to find the area of this parallelogram. Form vectors along these sides: F = i + j + 4k Calculate

G = i + 3j + 3k

 j k  1 4  = −9i + j + 2k 3 3  √ and the area of the parallelogram is F × G = 86 square units.   i  F × G =  1  1

If a rectangular box is skewed as in Figure 6.27, the resulting solid is called a rectangular parallelopiped. All of its faces are parallelograms. We can find the volume of such a solid by combining dot and cross products, as follows. F×G

ψ

H G F

FIGURE 6.27

FIGURE 6.28 Volume

Parallelopiped.

= H·
F×G.

THEOREM 6.7

Let F, G and H be vectors along incident sides of a rectangular parallelopiped. Then the volume of the parallelopiped is H ·
F × G. This is the absolute value of the real number formed by taking the dot product of H with F × G. Proof Figure 6.28 shows the parallelopiped. F × G is normal to the plane of F and G and oriented as shown in the diagram, according to the right-hand rule. If H is along the third side of the parallelogram, and  is the angle between H and F × G, then H cos
 is the altitude of the parallelopiped. The area of the base parallelogram is F × G by Theorem 6.6. Thus the volume of the parallelopiped is

H F × G cos
 But this is H ·
F × G.

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EXAMPLE 6.13

One corner of a rectangular parallelopiped is at
−1 2 2, and three incident sides extend from this point to
0 1 1,
−4 6 8, and
−3 −2 4. To find the volume of this solid, form the vectors F =
0 −
−1i +
1 − 2j +
1 − 2k = i − j − k G =
−4 −
−1i +
6 − 2j +
8 − 2k = −3i + 4j + 6k and H =
−3 −
−1i +
−2 − 2j +
4 − 2k = −2i − 4j + 2k Calculate

   i j k   F × G =  1 −1 −1  = −2i − 3j + k  −3 4 6 

Then H ·
F × G =
−2
−2 +
−4
−3 +
2
1 = 18 and the volume is 18 cubic units. The quantity H ·
F × G is called a scalar triple product. We will outline one of its properties in the problems.

SECTION 6.3

PROBLEMS

In each of Problems 1 through 6, compute F × G and, independently, G × F, verifying that one is the negative of the other. Use the dot product to compute the cosine of the angle  between F and G and use this to determine sin
. Then calculate F G sin
 and verify that this gives F × G.

11.
−4 2 −6
1 1 3
−2 4 5 In each of Problems 12 through 16, find the area of the parallelogram having incident sides extending from the first point to each of the other two.

1. F = −3i + 6j + k G = −i − 2j + k

12.
1 −3 7
2 1 1
6 −1 2

2. F = 6i − k G = j + 2k

13.
6 1 1
7 −2 4
8 −4 3

3. F = 2i − 3j + 4k G = −3i + 2j

14.
−2 1 6
2 1 −7
4 1 1

4. F = 8i + 6j G = 14j

15.
4 2 −3
6 2 −1
2 −6 4

5. F = 5i + 3j + 4k G = 20i + 6k

16.
1 1 −8
9 −3 0
−2 5 2

6. F = 2k G = 8i − j

In each of Problems 17 through 21, find the volume of the parallepiped whose incident sides extend from the first point to each of the other three.

In each of Problems 7 through 11, determine whether the points are collinear. If they are not, find an equation of the plane containing all three points.

17.
1 1 1
−4 2 7
3 5 7
0 1 6

7.
−1 1 6
2 0 1
3 0 0

18.
0 1 −6
−3 1 4
1 7 2
−3 0 4

8.
4 1 1
−2 −2 3
6 0 1

19.
1 6 1
−2 4 2
3 0 0
2 2 −4

9.
1 0 −2
0 0 0
5 1 1

20.
0 1 7
9 1 3
−2 4 1
3 0 −3

10.
0 0 2
−4 1 0
2 −1 1

21.
1 1 1
2 2 2
6 1 3
−2 4 6

6.4 The Vector Space R n In each of Problems 22 through 26, find a vector normal to the given plane. There are infinitely many such vectors. 22. 8x − y + z = 12 23. x − y + 2z = 0

30. Use vector operations to find a formula for the area of the triangle having vertices
ai  bi  ci  for i = 1 2 3. What conditions must be placed on these coordinates to ensure that the points are not collinear (all on a line)? The scalar triple product of F, G, and H is defined to be F G H = F ·
G × H.

24. x − 3y + 2z = 9 25. 7x + y − 7z = 7 26. 4x + 6y + 4z = −5 27. Prove that F ×
G + H = F × G + F × H. 28. Prove that
F × G = F ×
G = 
F × G. 29. Prove that F ×
G × H =
F · HG −
F · GH.

6.4

223

31. Let F = a1 i + b1 j + c1 k, G = a2 i + b2 j + c2 k, and H = a3 i + b3 j + c3 k. Prove that    a1 b1 c1   

F G H =  a2 b2 c2    a3 b3 c3 

The Vector Space Rn The world of everyday experience has three space dimensions. But often we encounter settings in which more dimensions occur. If we want to specify not only the location of a particle but the time in which it occupies a particular point, we need four coordinates
x y z t. And specifying the location of each particle in a system of particles may require any number of coordinates. The natural setting for such problems is R n , the space of points having n coordinates.

DEFINITION 6.8

n-vector

If n is a positive integer, an n-vector is an n-tuple
x1  x2      xn , with each coordinate xj a real number. The set of all n-vectors is denoted R n .

R1 is the real line, consisting of all real numbers. We can think of real numbers as 1-vectors, but there is is no advantage to doing this. R 2 consists of ordered pairs
x y of real numbers, and each such ordered pair (or 2-vector) can be identified with a point in the plane. R 3 consists of all 3-vectors, or points in 3-space. If n ≥ 4, we can no longer draw a set of mutually independent coordinate axes, one for each coordinate, but we can still work with vectors in R n according to rules we will now describe.

DEFINITION 6.9

Algebra of R n

1. Two n-vectors are added by adding their respective components:
x1  x2      xn  +
y1  y2      yn  =
x1 + y1  x2 + y2      xn + yn  2. An n-vector is multiplied by a scalar by multiplying each component by the scalar: 
x1  x2      xn  =
x1  x2      xn 

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The zero vector in R n is the n-vector O =
0 0     0 having each coordinate equal to zero. The negative of F =
x1  x2      xn  is −F =
−x1  −x2      −xn . As we did with n = 3, we denote G +
−F as G − F. The algebraic rules in R n mirror those we saw for R 3 . THEOREM 6.8

Let F, G, and H be in R n , and let  and  be real numbers. Then 1. F + G = G + F. 2. F +
G + H =
F + G + H. 3. F + O = F. 4. 5. 6. 7.


 + F = F + F.
F = 
F. 
F + G = F + G. O = O.

Because of these properties of the operations of addition of n-vectors, and multiplication of an n-vector by a scalar, we call R n a vector space. In the next section we will clarify the sense in which R n can be said to have dimension n. The length (norm, magnitude) of F =
x1  x2      xn  is defined by a direct generalization from the plane and 3-space: " F = x12 + x22 + · · · + xn2  There is no analogue of the cross product for vectors in R n when n > 3. However, the dot product readily extends to n-vectors.

DEFINITION 6.10

Dot Product of n-Vectors

The dot product of
x1  x2      xn  and
y1  y2      yn  is defined by
x1  x2      xn  ·
y1  y2      yn  = x1 y1 + x2 y2 + · · · + xn yn 

All of the conclusions of Theorem 6.3 remain true for n-vectors, as does Lemma 6.1. We will record these results for completeness. THEOREM 6.9

Let F, G, and H be n-vectors, and let  and  be real numbers. Then 1. 2. 3. 4.

F · G = G · F.
F + G · H = F · H + G · H. 
F · G =
F · G = F ·
G. F · F = F2 .

6.4 The Vector Space R n

225

5. F · F = 0 if and only if F = O. 6. F + G2 = 2 F2 + 2F · G + 2 G2 . The Cauchy-Schwarz inequality holds for n-vectors, but the proof given previously for 3-vectors does not generalize to R n . We will therefore give a proof that is valid for any n. THEOREM 6.10

Cauchy-Schwarz Inequality in R n

Let F and G be in R n . Then F · G ≤ F G Proof The inequality reduces to 0 ≤ 0 if either vector is the zero vector. Thus suppose F = O and G = O. Choose  = G and  = − F in Theorem 6.9(6). We get

0 ≤ F + G2 = G2 F2 − 2 G F F · G + F2 G2  Upon dividing this inequality by 2 F G we obtain F · G ≤ F G  Now go back to Theorem 6.9(6), but this time choose  = G and  = F to get 0 ≤ G2 F2 + 2 G F F · G + G2 F2 and upon dividing by 2 F G we get − F G ≤ F · G We have now shown that − F G ≤ F · G ≤ F G  and this is equivalent to the Cauchy-Schwarz inequality. In view of the Cauchy-Schwarz inequality, we can define the cosine of the angle between vectors F and G in R n by  0 if F or G equals the zero vector cos
 = 
F · G/
F G if F = O and G = O This is sometimes useful in bringing some geometric intuition to R n . For example, it is natural to define F and G to be orthogonal if the angle between them is /2, and by this definition of cos
, this is equivalent to requiring that F · G = 0, consistent with orthogonality in R 2 and R 3 . We can define a standard representation of vectors in R n by defining unit vectors along the n directions: e1 =
1 0 0     0 e2 =
0 1 0     0   en =
0 0     0 1

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Now any n-vector can be written
x1  x2      xn  = x1 e1 + x2 e2 + · · · + xn en =

n 

xj e j 

j=1

A set of n-vectors containing the zero vector, as well as sums of vectors in the set and scalar multiples of vectors in the set, is called a subspace of R n .

DEFINITION 6.11

Subspace

A set S of n-vectors is a subspace of R n if: 1. O is in S. 2. The sum of any vectors in S is in S. 3. The product of any vector in S with any real number is also in S.

Conditions (2) and (3) of the definition can be combined by requiring that F + G be in S for any vectors F and G in S, and any real numbers  and .

EXAMPLE 6.14

Let S consist of all vectors in R n having norm 1. In R 2 (the plane) this is the set of points on the unit circle about the origin; in R 3 this is the set of points on the unit sphere about the origin. S is not a subspace of R n for several reasons. First, O is not in S because O does not have norm 1. Further, a sum of vectors in S is not in S (a sum of vectors having norm 1 does not have norm 1). And, if  = 1, and F has norm 1, then F does not have norm 1, so F is not in S. In this example S failed all three criteria for being a subspace. It is enough to fail one to disqualify a set of vectors from being a subspace.

EXAMPLE 6.15

Let K consist of all scalar multiples of
−1 4 2 0 in R 4 . We want to know if K is a subspace of R 4 . First, O is in K, because O = 0
−1 4 2 0 =
0 0 0 0. Next, if F and G are in K, then F = 
−1 4 2 0 for some  and G = 
−1 4 2 0 for some , so F + G =
 + 
−1 4 2 0 is a scalar multiple of
−1 4 2 0 and therefore is in K. Finally, if F = 
−1 4 2 0 is any vector in K, and  is any scalar, then F =

−1 4 2 0 is a scalar multiple of
−1 4 2 0, and hence is in K. Thus K is a subspace of R 4 .

6.4 The Vector Space R n

227

EXAMPLE 6.16

Let S consist of just the zero vector O in R n . Then S is a subspace of R n . This is called the trivial subspace. At the other extreme, R n is also a subspace of R n . In R 2 and R 3 there are simple geometric characterizations of all possible subspaces. Begin with R 2 . We claim that only subspaces are the trivial subspace consisting of just the zero vector, or R 2 itself, or all vectors lying along a single line through the origin. To demonstrate this, we need the following fact. LEMMA 6.2

Let F and G be nonzero vectors in R 2 that are not parallel. Then every vector in R 2 can be written in the form F + G for some scalars  and . Represent F and G as arrows from the origin (Figure 6.29). These determine nonparallel lines L1 and L2 , respectively, through the origin, because F and G are assumed to be nonparallel. Let V be any 2-vector. If V = O, then V = 0F + 0G. We therefore consider the case that V = O and represent it as an arrow from the origin as well. We want to show that V must be the sum of scalar multiples of F and G. If V is along L1 , then V = F for some real number , and then V = F + 0G. Similarly, if V is along L2 , then V = G = 0F + G. Thus, suppose that V is not a scalar multiple of either F or G. Then the arrow representing V is not along L1 or L2 . Now carry out the construction shown in Figure 6.30. Draw lines parallel to L1 and L2 from the tip of V. Arrows from the origin to where these parallels intersect L1 and L2 determine, respectively, vectors A and B. By the parallelogram law, Proof

V = A + B But A is along L1 , so A = F for some scalar . And B is along L2 , and so B = G for some scalar . Thus V = F + G, completing the proof. We can now completely characterize the subspaces of R 2 . y B y

Parallel to L 1

L1 F x

V F

G A

L2

x G FIGURE 6.29

THEOREM 6.11

Parallel to L2 FIGURE 6.30

The Subspaces of R 2

Let S be a subspace of R 2 . Then one of the following three possibilities must hold: 1. S = R2 , or 2. S consists of just the zero vector, or 3. S consists of all vectors parallel to some straight line through the origin.

CHAPTER 6

228

Vectors and Vector Spaces

Suppose cases (1) and (2) do not hold. Because S is not the trivial subspace, S must contain at least one nonzero vector F. We will show that every vector in S is a scalar multiple of F. Suppose instead that there is a nonzero vector G in S that is not a scalar multiple of F. If V is any vector in R 2 , then by Lemma 6.2, V = F + G for some scalars  and . But S is a subspace or R 2 , so F + G is in S. This would imply that every vector in R 2 is in S, hence that S = R2 , a contradiction. Therefore there is no vector in S that is not a scalar multiple of F. We conclude that every vector in S is a scalar multiple of F, and the proof is complete. Proof

By a similar argument involving more cases, we can prove that any subspace of R 3 must be either the trivial subspace, R 3 itself, all vectors parallel to some line through the origin, or all vectors parallel to some plane through the origin.

SECTION 6.4

PROBLEMS

In each of Problems 1 through 6, find the sum of the vectors and express this sum in standard form. Calculate the dot product of the vectors and the angle between them. The latter may be expressed as an inverse cosine of a number. 1.
−1 6 2 4 0
6 −1 4 1 1 2.
0 1 4 −3
2 8 6 −4 3.
1 −4 3 2
16 0 0 4 4.
6 1 1 −1 2
4 3 5 1 −2 5.
0 1 6 −4 1 2 9 −3
6 6 −12 4 −3 −3 2 7 6.
−5 2 2 −7 −8
1 1 1 −8 7 In each of Problems 7 through 13, determine whether the set of vectors is a subspace of R n for the appropriate n. 7. S consists of all vectors
x y z x x in R5 . 8. S consists of all vectors
x 2x 3x y in R 4 . 9. S consists of all vectors
x 0 0 1 0 y in R6 . 10. S consists of all vectors
0 x y in R 3 .

6.5

11. S consists of all vectors
x y x + y x − y in R 4 . 12. S consists of all vectors in R7 having zero third and fifth components. 13. S consists of all vectors in R 4 whose first and second components are equal. 14. Let S consist of all vectors in R 3 on or parallel to the plane ax + by + cz = k, with a, b, and c real numbers, at least one of which is nonzero, and k = 0. Is S a subspace of R 3 ? 15. Let F and G be in R n . Prove that  F + G2 + F − G2 = 2 F2 + G2  Hint: Use the fact that the square of the norm of a vector is the dot product of the vector with itself. 16. Let F and G be orthogonal vectors in R n . Prove that F + G2 = F2 + G2 . This is called Pythagoras’s theorem. 17. Suppose F and G are vectors in R n satisfying the relationship of Pythagoras’s theorem (Problem 16). Does it follow that F and G are orthogonal?

Linear Independence, Spanning Sets, and Dimension in Rn In solving systems of linear algebraic equations and systems of linear differential equations, as well as for later work in Fourier analysis, we will use terminology and ideas from linear algebra. We will define these terms in R n , where we have some geometric intuition.

6.5 Linear Independence, Spanning Sets, and Dimension in R n

DEFINITION 6.12

229

Linear Combinations in R n

A linear combination of k vectors F1      Fk in R n is a sum 1 F1 + · · · + k Fk in which each j is a real number.

For example, −8
−2 4 1 0 + 6
1 1 −1 7 − 
8 0 0 0 is a linear combination of
−2 4 1 0,
1 1 −1 7, and
8 0 0 0 in R 4 . This linear combination is equal to the 4-vector
22 − 8 −26 −14 42 The set of all linear combinations of any given (finite) number of vectors in R n is always a subspace of R n . THEOREM 6.12

Let F1      Fk be in R n , and let V consist of all vectors 1 F1 + · · · + k Fk , in which each j can be any real number. Then V is a subspace of R n . First, O is in V (choose 1 = 2 = · · · = k = 0). Next, suppose G and H are in V . Then

Proof

G = 1 F1 + · · · + k Fk and H = 1 F1 + · · · + k Fk for some real numbers 1      k  1      k . Then G + H =
1 + 1 F1 + · · · +
k + k Fk is again a linear combination of F1      Fk , and so is in V . Finally, let G = 1 F1 + · · · + k Fk be in V . If c is any real number, then cG =
c1 F1 + · · · +
ck Fk is also a linear combination of F1      Fk , and is therefore in V . Therefore V is a subspace of R n .

Whenever we form a subspace by taking a linear combination of given vectors, we say that these vectors span the subspace.

DEFINITION 6.13

Spanning Set

Let F1      Fk be vectors in a subspace S of R n . Then F1      Fk form a spanning set for S if every vector in S is a linear combination of F1      Fk . In this case we say that S is spanned by F1      Fk , or that F1      Fk span S.

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For example, i, j, and k span R 3 , because every vector in R 3 can be written ai + bj + ck. The vector i + j in R 2 spans the subspace consisting of all vectors 
i + j, with  any scalar. These vectors all lie along the straight line y = x through the origin in the plane. Different sets of vectors may span the same subspace of R n . Consider the following example.

EXAMPLE 6.17

Let V be the subspace of R 4 consisting of all vectors
  0 0. Every vector in S can be written
  0 0 = 
1 0 0 0 + 
0 1 0 0 so
1 0 0 0 and
0 1 0 0 span V . But we can also write any vector in V as  
2 0 0 0 +
0  0 0 2  so the vectors
2 0 0 0 and
0  0 0 also span V . Different numbers of vectors may also span the same subspace. The vectors
  0 0 =


4 0 0 0
0 3 0 0
1 2 0 0 also span V . To see this, write an arbitrary vector in V as
  0 0 =

−2 −4
4 0 0 0 +
0 3 0 0 + 2
1 2 0 0 4 3

The last example suggests that some spanning sets are more efficient than others. If two vectors will span a subspace V , why should we choose a spanning set with three vectors in it? Indeed, in the last example, the last vector in the spanning set
4 0 0 0,
0 3 0 0,
1 2 0 0 is a linear combination of the first two: 2 1
1 2 0 0 =
4 0 0 0 +
0 3 0 0 4 3 Thus any linear combination of these three vectors can always be written as a linear combination of just
4 0 0 0 and
0 3 0 0. The third vector, being a linear combination of the first two, is “extraneous information.” These ideas suggest the following definition.

DEFINITION 6.14

Linear Dependence and Independence

Let F1      Fk be vectors in R n . 1. F1      Fk are linearly dependent if and only if one of these vectors is a linear combination of the others. 2. F1      Fk are linearly independent if and only if they are not linearly dependent.

Linear dependence of F1      Fk means that, whatever information these vectors carry, not all of them are needed, because at least one of them can be written in terms of the others. For example, if Fk = 1 F1 + · · · + k−1 Fk−1 

6.5 Linear Independence, Spanning Sets, and Dimension in R n

231

then knowing just F1      Fk−1 gives us Fk as well. In this sense, linearly dependent vectors are redundant. We can remove at least one Fj and the remaining k − 1 vectors will span the same subspace of R n that F1      Fk do. Linear independence means that no one of the vectors F1      Fk is a linear combination of the others. Whatever these vectors are telling us (for example, specifying a subspace), we need all of them or we lose information. If we omit Fk , we cannot retrieve it from F1      Fk−1 .

EXAMPLE 6.18

The vectors
1 1 0 and
−2 0 3 are linearly independent in R 3 . To prove this, suppose instead that these vectors are linearly dependent. Then one is a linear combination of the other, say
−2 0 3 = 
1 1 0 But then, from the first components,  = −2, while from the second components,  = 0, an impossibility. These vectors span the subspace V of R 3 consisting of all vectors 
−2 0 3 + 
1 1 0 Both of the vectors
1 1 0 and
−2 0 3 are needed to describe V . If we omit one, say
1 1 0, then the subspace of R 3 spanned by the remaining vector,
−2 0 3, is different from V . For example, it does not have
1 1 0 in it. The following is a useful characterization of linear dependence and independence. THEOREM 6.13

Let F1      Fk be vectors in R n . Then 1. F1      Fk are linearly dependent if and only if there are real numbers 1      k , not all zero, such that 1 F1 + 2 F2 + · · · + k Fk = O 2. F1      Fk are linearly independent if and only if an equation 1 F1 + 2 F2 + · · · + k Fk = O can hold only if 1 = 2 = · · · = k = 0. To prove (1), suppose first that F1      Fk are linearly dependent. Then at least one of these vectors is a linear combination of the others. Say, to be specific, that

Proof

F1 = 2 F2 + · · · + k Fk  Then F1 − 2 F2 + · · · + k Fk = O But this is a linear combination of F1      Fk adding up to the zero vector and having a nonzero coefficient (the coefficient of F1 is 1). Conversely, suppose 1 F1 + 2 F2 + · · · + k Fk = O

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with at least some j = 0. We want to show that F1      Fk are linearly dependent. By renaming the vectors if necessary, we may suppose for convenience that 1 = 0. But then F1 = −

2  F2 − · · · − k Fk  1 1

so F1 is a linear combination of F2      Fk and hence F1      Fk are linearly dependent. This completes the proof of (1). Conclusion (2) follows from (1) and the fact that F1      Fk are linearly independent exactly when these vectors are not linearly dependent. This theorem suggests a strategy for determining whether a given set of vectors is linearly dependent or independent. Given F1      Fk , set 1 F1 + 2 F2 + · · · + k Fk = O

(6.4)

and attempt to solve for the coefficients 1      k . If equation (6.4) forces 1 = · · · = k = 0, then F1      Fk are linearly independent. If we can find at least one nonzero j so that equation (6.4) is true, then F1      Fk are linearly dependent.

EXAMPLE 6.19

Consider
1 0 3 1,
0 1 −6 −1 and
0 2 1 0 in R 4 . We want to know whether these vectors are linearly dependent or independent. Look at a linear combination c1
1 0 3 1 + c2
0 1 −6 −1 + c3
0 2 1 0 =
0 0 0 0 If this is to hold, then each component of the vector
c1  c2 + 2c3  3c1 − 6c2 + c3  c1 − c2  must be zero: c1 = 0 c2 + 2c3 = 0 3c1 − 6c2 + c3 = 0 c1 − c2 = 0 The first equation gives c1 = 0, so the fourth equation tells us that c2 = 0. But then the second equation requires that c3 = 0. Therefore, the only linear combination of these three vectors that equals the zero vector is the trivial linear combination (all coefficients zero), and by (2) of Theorem 6.13, the vectors are linearly independent. In the plane R 2 , two vectors are linearly dependent if and only if they are parallel. In R 3 , two vectors are linearly dependent if and only if they are parallel, and three vectors are linearly dependent if and only if they are in the same plane. Any set of vectors that includes the zero vector must be linearly dependent. Consider, for example, the vectors O, F2 ,   , Fk . Then the linear combination 1O + 0F2 + · · · + 0Fk = O is a linear combination of these vectors that add up to the zero vector, but has a nonzero coefficient (the coefficient of O is 1). By Theorem 6.13(1), these vectors are linearly dependent. There is a special circumstance in which it is particularly easy to tell that a set of vectors is linearly independent. This is given in the following lemma, which we will use later.

6.5 Linear Independence, Spanning Sets, and Dimension in R n

233

LEMMA 6.3

Let F1      Fk be vectors in R n . Suppose each Fj has a nonzero element in some component where each of the other Fi
s has a zero component. Then F1      Fk are linearly independent. An example will clarify why this is true.

EXAMPLE 6.20

Consider the vectors F1 =
0 4 0 0 2 F2 =
0 0 6 0 −5 F3 =
0 0 0 −4 12 To see why these are linearly independent, suppose F1 + F2 + !F3 =
0 0 0 0 0 Then
0 4 6 −4! 2 − 5 + 12! =
0 0 0 0 0 From the second components, 4 = 0 so  = 0. From the third components, 6 = 0 so  = 0. And from the fourth components, −4! = 0 so ! = 0. Then the vectors are linearly independent by Theorem 6.13(2). The fact that each of the vectors has a nonzero element where all the others have only zero components makes it particularly easy to conclude that  =  = ! = 0, and that is what is needed to apply Theorem 6.13. There is another important setting in which it is easy to tell that vectors are linearly independent. Nonzero vectors F1      Fk in R n are said to be mutually orthogonal if each is orthogonal to each of the other vectors in the set. That is, Fi ·Fj = 0 if i = j. Mutually orthogonal nonzero vectors are necessarily linearly independent.

THEOREM 6.14

Let F1      Fk be mutually orthogonal nonzero vectors in R n . Then F1      Fk are linearly independent. Proof

Suppose 1 F1 + 2 F2 + · · · + k Fk = O

For any j = 1     k,
1 F1 + 2 F2 + · · · + k Fk  · Fj = 0 = 1 F1 · Fj + 2 F2 · Fj + · · · + j Fj · Fj + · · · + k Fk · Fj # #2 = cj Fj · Fj = cj #Fj #  # #2 because Fi · Fj = 0 if i = j. But Fj is not the zero vector, so #Fj # = 0, hence cj = 0. Therefore, each coefficient is zero and F1      Fk are linearly independent by Theorem 6.13(2).

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EXAMPLE 6.21

The vectors
−4 0 0,
0 −2 1,
0 1 −2 are linearly independent in R 3 , because each is orthogonal to the other two. A “smallest” spanning set for a subspace of R n is called a basis for that subspace.

DEFINITION 6.15

Basis

Let V be a subspace of R n . A set of vectors F1      Fk in V form a basis for V if F1      Fk are linearly independent and also span V .

Thus, for F1      Fk to be a basis for V , every vector in V must be a linear combination of F1      Fk , and if any Fj is omitted from the list F1      Fk , the remaining vectors do not span V . In particular, if Fj is omitted, then the subspace spanned by F1      Fj−1  Fj+1     Fk cannot contain Fj , because by linear independence, Fj is not a linear combination of F1      Fj−1  Fj+1     Fk .

EXAMPLE 6.22

i, j, and k form a basis for R 3 , and e1  e2      en form a basis for, R n .

EXAMPLE 6.23

Let V be the subspace of R n consisting of all n−vectors with zero first component. Then e2      en form a basis for V .

EXAMPLE 6.24

In R 2 , let V consist of all vectors parallel to the line y = 4x. Every vector in V is a multiple of
1 4. This vector by itself forms a basis for V . In fact, any vector
 4 with  = 0 forms a basis for V .

EXAMPLE 6.25

In R 3 , let M be the subspace of all vectors on or parallel to the plane x + y + z = 0. A vector
x y z in R 3 is in M exactly when z = −x − y, so such a vector can be written
x y z =
x y −x − y = x
1 0 −1 + z
0 1 −1 The vectors
1 0 −1 and
0 1 −1 therefore span M. Since these two vectors are linearly independent, they form a basis for M. We may think of a basis of V as a minimal linearly independent spanning set F1      Fk for V . If we omit any of these vectors, the remaining vectors will not be enough to span V .

6.5 Linear Independence, Spanning Sets, and Dimension in R n

235

And if we use additional vectors, say the set F1      Fk  H, then this set also spans V , but is not linearly independent (because H is a linear combination of F1      Fk ). There is nothing unique about a basis for a subspace of R n . Any nontrivial subspace of R n has infinitely many different bases. However, it is a theorem of linear algebra, which we will not prove, that for a given subspace V of R n , every basis has the same number of vectors in it. This number is the dimension of the subspace.

DEFINITION 6.16

The dimension of a subspace of R n is the number of vectors in any basis for the subspace.

In particular, R n (which is a subspace of itself) has dimension n, a basis consisting of the n vectors e1      en . The subspace in Example 6.25 has dimension 2.

SECTION 6.5

PROBLEMS

In each of Problems 1 through 10, determine whether the given vectors are linearly independent or dependent in R n for appropriate n. 1. 3i + 2j i − j in R

13. i + 6j − 2k −i + 4j − 3k i + 16j − 7k 14. 4i − 3j + k 10i − 3j 2i − 6j + 3k

3

2. 2i 3j 5i − 12k i + j + k in R 3 3.
8 0 2 0 0 0 0
0 0 0 0 1 −1 0 in R7 4.
1 0 0 0
0 1 1 0
−4 6 6 0 in R 4 5.
1 2 −3 1,
4 0 0 2,
6 4 −6 4 in R 4 6.
0 1 1 1
−3 2 4 4
−2 2 34 2
1 1 −6 −2 in R 4 7.
1 −2
4 1
6 6 in R 2

15. 8i + 6j 2i − 4j i + k 16. 12i − 3k i + 2j − k −3i + 4j In each of Problems 17 through 24, determine a basis for the subspace S of R n and determine the dimension of the subspace. 17. S consists of all vectors
x y −y −x in R 4 . 18. S consists of all vectors
x y 2x3y in R 4 .

8.
−1 1 0 0 0
0 −1 1 0 0
0 1 1 1 0 in R5

19. S consists of all vectors in the plane 2x − y + z = 0.

9.
−2 0 0 1 1
1 0 0 0 0
0 0 0 0 2
1 −1 3 3 1 in R 4

20. S consists of all vectors
x y −y x − y z in R5 .

10.
3 0 0 4
2 0 0 8 in R 4

21. S consists of all vectors in R 4 with zero second component.

11. Prove that three vectors in R 3 are linearly dependent if and only if their scalar triple product is zero. (See Problem 31 in Section 6.3).

22. S consists of all vectors
−x x y 2y in R 4 .

In each of Problems 12 through 16, use the result of Problem 11 to determine whether the three vectors in R 3 are linearly dependent or independent. 12. 3i + 6j − k 8i + 2j − 4k i − j + k

23. S consists of all vectors parallel to the line y = 4x in R 2 . 24. S consists of all vectors parallel to the plane 4x + 2y − z = 0 in R 3 .

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CHAPTER

MATRICES ELEMENTARY ROW OPERATIONS AND ELEMENTARY MATRICES THE ROW ECHELON FORM OF A MATRIX SOLUTION OF HOMOGENEOUS SYSTEMS OF LINEAR EQUATIONS THE SOLUTION SPACE OF

7

Matrices and Systems of Linear Equations

This chapter is devoted to the notation and algebra of matrices, as well as their use in solving systems of linear algebraic equations. To illustrate the idea of a matrix, consider a system of linear equations: x1 + 2x2 − x3 + 4x4 = 0 3x1 − 4x2 + 2x3 − 6x4 = 0 x1 − 3x2 − 2x3 + x4 = 0 All of the information needed to solve this system lies in its coefficients. Whether the first unknown is called x1 , or y1 , or some other name, is unimportant. It is important, however, that the coefficient of the first unknown in the second equation is 3. If we change this number we may change the solutions of the system. We can therefore work with such a system by storing its coefficients in an array called a matrix: ⎞ ⎛ 1 2 −1 4 ⎝3 −4 2 −6⎠  1 −3 −2 1 This matrix displays the coefficients in the pattern in which they appear in the system of equations. The coefficients of the ith equation are in row i, and the coefficients of the j th unknown xj are in column j. The number in row i, column j is the coefficient of xj in equation i. But matrices provide more than a visual aid or storage device. The algebra and calculus of matrices will form the basis for methods of solving systems of linear algebraic equations, and later for solving systems of linear differential equations and analyzing solutions of systems of nonlinear differential equations.

237

238

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Matrices and Systems of Linear Equations

Matrices DEFINITION 7.1

Matrix

An n by m matrix is an array of objects arranged in n rows and m columns.

We will denote matrices by boldface type, as was done with vectors. When A is an n by m matrix, we often write that A is n × m (read “n by m”). The first integer is the number of rows in the matrix, and the second integer is the number of columns. The objects in the matrix may be numbers, functions, or other quantities. For example,   2 √1  2 −5 1 is a 2 × 3 matrix,



e2x cos
x

is a 2 × 2 matrix, and

e−4x x2



⎛

⎞ 0 ⎜−4 ⎟ ⎜ 3⎟ ⎝ x⎠ 2

is a 4 × 1 matrix. A matrix having the same number of rows as columns is called a square matrix. The 2 × 2 matrix shown above is square. The object in row i and column j of a matrix is called the i j element, or i j entry of the matrix. If a matrix is denoted by an upper case letter, say A, then its i j element is often denoted aij and we write A = aij . For example, if ⎛ ⎞ 0 x H = hij  = ⎝1 − sin
x 1 − 2i⎠ x2 i then H is a 3 × 2 matrix, and h11 = 0, h12 = x, h21 = 1 − sin
x, h22 = 1 − 2i, h31 = x2 and h32 = i. We will be dealing with matrices whose elements are real or complex numbers, or functions. Sometimes it is also convenient to denote the i j element of A by
Aij . In the matrix H,
H22 = 1 − 2i DEFINITION 7.2

and
H31 = x2 

Equality Matrices

A = aij  and B = bij  are equal if and only if they have the same number of rows, the same number of columns, and for each i and j, aij = bij .

If two matrices either have different numbers of rows or columns, or if the objects in a particular location in the matrices are different, then the matrices are unequal.

7.1 Matrices

7.1.1

239

Matrix Algebra

We will develop the operations of addition and multiplication of matrices and multiplication of a matrix by a number.

DEFINITION 7.3

Matrix Addition

If A = aij  and B = bij  are n × m matrices, then their sum is the n × m matrix A + B = aij + bij 

We therefore add matrices by adding corresponding elements. For example,       1 2 −3 −1 6 3 0 8 0 + =  4 0 2 8 12 14 12 12 16 If two matrices are of different dimensions (different numbers of rows or columns), then they cannot be added, just as we do not add 4-vectors and 7-vectors.

DEFINITION 7.4

Product of a Matrix and a Scalar

If A = aij  and  is a scalar, then A is the matrix defined by A = aij 

This means that we multiply a matrix by  by multiplying each element of the matrix by . For example, ⎛ ⎞ ⎛ ⎞ 2 0 6 0 ⎜0 0⎟ ⎜0 0 ⎟ ⎟ ⎜ ⎟ 3⎜ ⎝1 4⎠ = ⎝3 12⎠ 2 6 6 18 and

 x

  1 x x = −x2 −x cos
x

 x2  x cos
x

Some, but not all, pairs of matrices can be multiplied.

DEFINITION 7.5

Multiplication of Matrices

Let A = aij  be an n × r matrix, and B = bij  an r × m matrix. Then the matrix product AB is the n × m matrix whose i j element is ai1 b1j + ai2 b2j + · · · + air brj  That is,

 AB =

r 

k=1

 aik bkj 

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If we think of each row of A as an r-vector, and each column of B as an r-vector, then the i j element of AB is the dot product of row i of A with column j of B : i j element of AB =
row i of A ·
column j of B This is why the number of columns of A must equal the number of rows of B for AB to be defined. These rows of A and columns of B must be vectors of the same length in order to take this dot product. Thus not every pair of matrices can be multiplied. Further, even when AB is defined, BA need not be. We will give one rationale for defining matrix multiplication in this way shortly. First we will look at some examples of matrix products and then develop the rules of matrix algebra.

EXAMPLE 7.1

Let

 A=

1 2

  1 3 and B = 5 2

1 1

 3  4

Then A is 2 × 2 and B is 2 × 3, so AB is defined (number of columns of A equals the number of rows of B). Further, AB is 2 × 3 (number of rows of A, number of columns of B). Now compute    1 3 1 1 3 AB = 2 5 2 1 4  
1 3 ·
1 2
1 3 ·
1 1
1 3 ·
3 4 =
2 5 ·
1 2
2 5 ·
1 1
2 5 · 3 4   7 4 15 =  12 7 26 In this example, BA is not defined, because the number of columns of B, which is 3, does not equal the number of rows of A, which is 2.

EXAMPLE 7.2

Let  A=

1 4

1 2 1 6

1 2



⎛

−1 ⎜ 2 and B = ⎜ ⎝ 1 12

⎞ 8 1⎟ ⎟ 1⎠ 6

Since A is 2 × 4 and B is 4 × 2, AB is defined and is a 2 × 2 matrix:  
1 1 2 1 ·
−1 2 1 12
1 1 2 1 ·
8 1 1 6 AB =
4 1 6 2 ·
−1 2 1 12
4 1 6 2 ·
8 1 1 6   15 17 =  28 51

7.1 Matrices In this example BA is also defined, and is ⎛ −1 ⎜ 2 BA = ⎜ ⎝ 1 12 ⎛ 31 ⎜ 6 =⎜ ⎝ 5 36

4 × 4: ⎞ 8  1⎟ ⎟ 1 1⎠ 4 6 7 46 3 10 2 8 18 60

1 2 1 6

1 2

241



⎞ 15 4⎟ ⎟ 3⎠ 24

As the last example shows, even when both AB and BA are defined, these may be matrices of different dimensions. Matrix multiplication is noncommutative, and it is the exception rather than the rule to have AB equal BA. If A is a square matrix, then AA is defined and is also square. Denote AA as A2 . Similarly, A
A2  = A3 and, for any positive integer k, Ak = AA · · · A, a product with k factors. Some of the rules for manipulating matrices are like those for real numbers.

THEOREM 7.1

Let A, B, and C be matrices. Then, whenever the indicated operations are defined, we have: 1. 2. 3. 4.

A + B = B + A. A
B + C = AB + AC.
A + BC = AC + BC. A
BC =
ABC.

For (1), both matrices must have the same dimensions, say n × m. For (2), B and C must have the same dimensions, and the number of columns in A must equal the number of rows in B and in C. For (4), A must be n × r, B must be r × k and C must be k × m. Then A
BC and
ABC are n × m. The theorem is proved by direct appeal to the definitions. We will provide the details for (1) and (2). To prove (1), let A = aij  and B = bij . Then

Proof

A + B = aij + bij  = bij + aij  = B + A because each aij and bij is a number or function and the addition of these objects is commutative. For (2), let A = aij , B = bij  and C = cij . Suppose A is n × k and B and C are k × m. Then B + C is k × m, so A
B + C is defined and is n × m. And AB and BC are both defined and n × m. There remains to show that the i j element of AB + AC is the same as the i j element of A
B + C. Row i of A, and columns j of B and C, are k-vectors, and from properties of the dot product, i j element of A
B + C =
row i of A ·
column j of B + C =
row i of A ·
column j of B +
row i of A ·
column j of C =
i j element of AB +
i j element of AC = i j element of AB + AC

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We have already noted that matrix multiplication does not behave in some ways like multiplication of numbers. Here is a summary of three significant differences. Difference 1 For matrices, even when AB and BA are both defined, possibly AB = BA.

EXAMPLE 7.3



but

1 −2



0 4



−2 1



−2 1

6 3

1 −2

  6 −2 = 3 8

  0 −14 = 4 −5

6 0



 24  12

Difference 2 There is no cancellation in products. If AB = AC, we cannot infer that B = C.

EXAMPLE 7.4





1 3

1 3

But

    2 1 1 2 7 = 16 3 3 5 11   7 18 =  21 54

4 3



  2 2  = 16 5

4 3

 7  11

Difference 3 The product of two nonzero matrices may be zero.

EXAMPLE 7.5



1 0

7.1.2

2 0



  6 4 0 = −3 −2 0

 0  0

Matrix Notation for Systems of Linear Equations

Matrix notation is very efficient for writing systems of linear algebraic equations. Consider, for example, the system 2x1 − x2 + 3x3 + x4 = 1 x1 + 3x2 − 2x4 = 0 −4x1 − x2 + 2x3 − 9x4 = −3 The matrix of coefficients of this system is the 3 × 4 matrix ⎞ ⎛ 2 −1 3 1 3 0 −2⎠  A=⎝ 1 −4 −1 2 −9

7.1 Matrices

243

Row i contains the coefficients of the ith equation, and column j contains the coefficients of xj . Define ⎛ ⎞ ⎛ ⎞ x1 1 ⎜ x2 ⎟ ⎟ and B = ⎝ 0⎠  X=⎜ ⎝ x3 ⎠ −3 x4 Then

⎛ ⎞ ⎞ x 1 ⎜ 1⎟ x2 ⎟ −2⎠ ⎜ ⎝ x3 ⎠ −9 x4 ⎛ ⎞ ⎛ ⎞ 2x1 − x2 + 3x3 + x4 1 ⎠ = ⎝ 0⎠  x1 + 3x2 − 2x4 =⎝ −3 −4x1 − x2 + 2x3 − 9x4 ⎛

2 AX = ⎝ 1 −4

−1 3 −1

3 0 2

We can therefore write the system of equations in matrix form as AX = B This is more than just notation. Soon this matrix formulation will enable us to use matrix operations to solve the system. A similar approach can be taken toward systems of linear differential equations. Consider the system x1
+ tx2
− x3
= 2t − 1 t2 x1
− cos
tx2
− x3
= et  Let 

1 A= 2 t

t − cos
t

 −1  −1

⎛ ⎞ x1 X = ⎝x2 ⎠  x3

  2t − 1 and F =  et

Then the system can be written AX
= F in which X
is formed by differentiating each matrix element of X. As with systems of linear algebraic equations, this formulation will enable us to bring matrix methods to bear on solving the system of differential equations. In both of these formulations, the definition of matrix product played a key role. Matrix multiplication may seem unmotivated at first, but it is just right for converting a system of linear algebraic or differential equations to a matrix equation.

7.1.3

Some Special Matrices

Some matrices occur often enough to warrant special names and notation.

DEFINITION 7.6

Zero Matrix

Onm denotes the n × m zero matrix, having each element equal to zero.

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For example,

 O23 =

0 0

0 0

 0  0

If A is n × m, then A + Onm = Onm + A = A The negative of A is the matrix obtained by replacing each element of A with its negative. This matrix is denoted −A. If A = aij , then −A = −aij . If A is n × m, then A +
−A = Onm  Usually we write A +
−B as A − B.

DEFINITION 7.7

Identity Matrix

The n × n identity matrix is the matrix In having each i j element equal to zero if i = j, and each i i− element equal to 1.

For example, 

1 I2 = 0

0 1

⎛



1 I 3 = ⎝0 0

and

0 1 0

⎞ 0 0⎠  1

THEOREM 7.2

If A is n × m, then AIm = In A = A We leave a proof of this to the student.

EXAMPLE 7.6

Let

⎛

1 A=⎝ 2 −1 Then

⎛

1 I 3 A = ⎝0 0 and

⎛

1 AI2 = ⎝ 2 −1

0 1 0

⎞⎛ 0 1 0⎠ ⎝ 2 1 −1 ⎞ 0  1 1⎠ 0 8

⎞ 0 1⎠  8 ⎞ ⎛ 0 1 1⎠ = ⎝ 2 8 −1



⎛

1 0 =⎝ 2 1 −1

⎞ 0 1⎠ = A 8

⎞ 0 1⎠ = A 8

7.1 Matrices

DEFINITION 7.8

245

Transpose

If A = aij  is an n × m matrix, then the transpose of A is the m × n matrix At = aji .

The transpose of A is formed by making row k of A, column k of At .

EXAMPLE 7.7

Let



−1 A= 0

6 

3 12

 3  −5

This is a 2 × 4 matrix. The transpose is the 4 × 2 matrix ⎞ ⎛ −1 0 ⎜ 6 ⎟ ⎟ At = ⎜ ⎝ 3 12⎠  3 −5

THEOREM 7.3

1.
In t = In . 2. For any matrix A,
At t = A. 3. If AB is defined, then
ABt = Bt At . Conclusion (1) should not be surprising, since row i of In is the same as column i, so interchanging rows and columns has no effect. Similarly, (2) is intutively clear. If we interchange rows and columns of A to form At , and then interchange the rows and columns of At , we should put everything back where it was, resulting in A again. We will prove conclusion (3). Proof Let A = aij  be n × k and let B = bij  be k × m. Then AB is defined and is n × m. Since Bt is m × k and At is k × n, then Bt At is defined and is m × n. Thus
ABt and Bt At have the same dimensions. Now we must show that the i j element of
ABt equals the i j element of Bt At . Falling back on the definition of matrix product, we have

i j element of Bt At =

k  s=1

=

k  s=1

This completes the proof of (3).


Bt is
At sj =

k 

bsi ajs

s=1

ajs bsi = j i element of AB = i j element of
ABt 

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In some calculations it is convenient to write the dot product of two n-vectors as a matrix product, using the transpose. Write the n-vector
x1  x2      xn  as a 1 × n column matrix: ⎛ ⎞ x1 ⎜ x2 ⎟ ⎜ ⎟ X = ⎜  ⎟ ⎝  ⎠ xn Then

 X t = x1

···

x2

 xn 

an n × 1 matrix. Let
y1  y2  · · ·  yn  also be an n-vector, which we write as an 1 × n column matrix: ⎛ ⎞ y1 ⎜ y2 ⎟ ⎜ ⎟ Y = ⎜  ⎟ ⎝  ⎠ yn Then Xt Y is the 1 × 1 matrix  X t Y = x1

x2

···

⎛ ⎞ y1 ⎜ y  ⎜ 2⎟ ⎟ xn ⎜  ⎟ ⎝  ⎠ yn

=
x1 y1 + x2 y2 + · · · + xn yn  = X · Y Here we have written the resulting 1 × 1 matrix as just its single element, without the matrix brackets. This is common practice for 1 × 1 matrices. We now have the dot product of two n-vectors, written as 1 × n column vectors, as the matrix product Xt Y This will prove particularly useful when we treat eigenvalues of matrices.

7.1.4

Another Rationale for the Definition of Matrix Multiplication

We have seen that matrix multiplication allows us to write linear systems of algebraic and differential equations in compact matrix form as AX = B or AX
= F. Matrix products are also tailored to other purposes, such as changes of variables in linear equations. To illustrate, consider a 2 × 2 linear system a11 x1 + a12 x2 = c1 a21 x1 + a22 x2 = c2 

(7.1)

Change variables by putting x1 = h11 y1 + h12 y2 x2 = h21 y1 + h22 y2  Then a11
h11 y1 + h12 y2  + a12
h21 y1 + h22 y2  = c1

(7.2)

7.1 Matrices

247

and a21
h11 y1 + h12 y2  + a22
h21 y1 + h22 y2  = c2  After rearranging terms, the transformed system is
a11 h11 + a12 h21 y1 +
a11 h12 + a12 h22 y2 = c1
a21 h11 + a22 h21 y1 +
a21 h12 + a22 h22 y2 = c2 

(7.3)

Now carrry out the same transformation using matrices. Write the original system (7.1) as AX = C, where       a a12 x c A = 11  X= 1 and C = 1 a21 a22 x2 c2 and the equations of the transformation (7.2) as X = HY, where     h h12 y H = 11 and Y = 1  h21 h22 y2 Then AX = A
HY =
AHY = C Now observe that



  a11 a12 h11 h12 AH = a21 a22 h21 h22   a11 h11 + a12 h21 a11 h12 + a12 h22 =  a21 h11 + a22 h21 a21 h12 + a22 h22

exactly as we found in the system (7.3) by term by term substitution. The definition of matrix product is just what is needed to carry out a linear change of variables. This idea also applies to linear transformations in systems of differential equations.

7.1.5

Random Walks in Crystals

We will conclude this section with another application of matrix multiplication, this time to the problem of enumerating the paths atoms can take through a crystal lattice. Crystals have sites arranged in a lattice pattern. An atom may jump from a site it occupies to any adjacent, vacant site. If more than one adjacent site is vacant, the atom “selects” its target site at random. The path such an atom makes through the crystal is called a random walk. We can represent the lattice of locations and adjacencies by drawing a point for each location, with a line between two points only if an atom can move directly from one to the other in the crystal. Such a diagram is called a graph. Figure 7.1 shows a typical graph. In this graph an atom could move from point v1 to v2 or v3 , to which it is connected by lines, but not directly to v6 because there is no line between v1 and v6 . Two points are called adjacent in G if there is a line between them in the graph. A point is not considered adjacent to itself—there are no lines starting and ending at the same point. A walk of length n in such a graph is a sequence t1      tn+1 of points (not necessarily different) with each tj adjacent to tj+1 in the graph. Such a walk represents a possible path an atom might take through various sites in the crystal. Points may repeat in a walk because an atom may return to the same site any number of times. A vi − vj walk is a walk that begins at vi and ends at vj .

248

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Matrices and Systems of Linear Equations v2 v1

v3

v6 v5

v4

FIGURE 7.1

A typical graph.

Physicists and materials engineers who study crystals are interested in the following question: given a crystal with n sites labeled v1      vn , how many different walks of length k are there between any two sites (or from a site back to itself)? Matrices enter into the solution of this problem as follows. Define the adjacency matrix A of the graph to be the n × n matrix having each i i element zero, and for i = j, the i j-element equal to 1 if there is a line in the graph between vi and vj , and 0 if there is no such line. For example, the graph of Figure 7.1 has adjacency matrix ⎞ ⎛ 0 1 1 1 0 0 ⎜1 0 1 0 0 0 ⎟ ⎟ ⎜ ⎜1 1 0 1 0 0 ⎟ ⎟ ⎜ A=⎜ ⎟ ⎜1 0 1 0 1 1 ⎟ ⎝0 0 0 1 0 1⎠ 0 0 0 1 1 0 The 1 2 element of A is 1 because there is a line between v1 and v2 , while the 1 5 element is zero because there is no line between v1 and v5 . The following remarkable theorem uses the adjacency matrix to solve the walk-enumeration problem.

THEOREM 7.4

Let A = aij  be the adjacency matrix of a graph G having points v1      vn . Let k be any positive integer. Then the number of distinct vi − vj walks of length k in G is equal to the i j element of Ak . We can therefore calculate the number of random walks of length k between any two points (or from any point back to itself) by reading the elements of the kth power of the adjacency matrix. Proof Proceed by mathematical induction on k. First consider the case k = 1. If i = j, there is a vi − vj walk of length 1 in G exactly when there is a line between vi and vj , and in this case aij = 1. There is no vi − vj walk of length 1 if vi and vj have no line between them, and in this case aij = 0. If i = j, there is no vi − vi walk of length 1, and aii = 0. Thus, in the case k = 1, the i j element of A gives the number of walks of length 1 from vi to vj , and the conclusion of the theorem is true. Now assume that the conclusion of the theorem is true for walks of length k. We will prove that the conclusion holds for walks of length k + 1. Thus, we are assuming that the i j

7.1 Matrices v3

v2

v8

v1 th

1

vr

g len

v7

length k

vi FIGURE 7.2

v4

v5

v6

vj

249

FIGURE 7.3

element of Ak is the number of distinct vi − vj walks of length k in G, and we want to prove that the i j element of Ak+1 is the number of distinct vi − vj walks of length k + 1. Consider how a vi − vj walk of length k + 1 is formed. First there must be a vi − vr walk of length 1 from vi to some point vr adjacent to vi , followed by a vr − vj walk of length k (Figure 7.2). Therefore number of distinct vi − vj walks of length k + 1 = sum of the number of distinct vr − vj walks of length k with the sum taken over all points vr adjacent to vi . Now air = 1 if vr is adjacent to vi , and air = 0 otherwise. Further, by the inductive hypothesis, the number of distinct vr − vj walks of length k is the r j element of Ak . Denote Ak = B = bij . Then, for r = 1    n, air brj = 0 if vr is not adjacent to vi and air brj = the number of distinct vi − vj walks of length k + 1 passing through vr , if vr is adjacent to vi . Therefore the number of vi − vj walks of length k + 1 in G is ai1 b1j + ai2 b2j + · · · + ain bnj because this counts the number of walks of length k from vr to vj for each point vr adjacent to vi . But this sum is exactly the i j element of AB, which is Ak+1 . This completes the proof by induction. For example, the adjacency matrix of the graph of Figure 7.3 is ⎛

0 ⎜1 ⎜ ⎜0 ⎜ ⎜0 A=⎜ ⎜0 ⎜ ⎜1 ⎜ ⎝0 0

1 0 1 0 0 0 1 1

0 1 0 1 0 0 0 0

0 0 1 0 1 1 1 1

0 0 0 1 0 1 1 0

1 0 0 1 1 0 0 0

0 1 0 1 1 0 0 1

⎞ 0 1⎟ ⎟ 0⎟ ⎟ 1⎟ ⎟ 0⎟ ⎟ 0⎟ ⎟ 1⎠ 0

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Suppose we want the number of v4 − v7 ⎛ 0 5 ⎜6 2 ⎜ ⎜1 7 ⎜ ⎜4 4 A3 = ⎜ ⎜2 5 ⎜ ⎜4 4 ⎜ ⎝3 9 2 8

walks of length 3 in G. Calculate ⎞ 1 4 2 4 3 2 7 4 5 4 9 8⎟ ⎟ 0 8 3 2 3 2⎟ ⎟ 8 6 8 8 11 10⎟ ⎟ 3 8 4 6 8 4⎟ ⎟ 2 8 6 2 4 4⎟ ⎟ 3 11 8 4 6 7⎠ 2 10 4 4 7 4

We read from the 4 7 element of A3 that there are 11 walks of length 3 from v4 to v7 . For this relatively simple graph, we can actually list all these walks: v4 v7 v4 v7  v4 v3 v4 v7  v4 v8 v4 v7  v4 v5 v4 v7  v4 v6 v4 v7  v4 v7 v8 v7  v4 v7 v5 v7  v4 v7 v2 v7  v4 v3 v2 v7  v4 v8 v2 v7  v4 v6 v5 v7  Obviously it would not be practical to determine the number of vi − vj walks of length k by explicitly listing them if k or n is large. Software routines for matrix calculations make this theorem a practical solution to the random walk counting problem.

SECTION 7.1

PROBLEMS

In each of Problems 1 through 6, carry out the computation with the given matrices A and B. ⎞ ⎛ ⎛ 1 −1 3 −4 0 ⎟ ⎜ ⎜ 6⎠  B = ⎝−2 −1 1. A = ⎝ 2 −4 −1 1 2 8 15 2A − 3B ⎞ ⎛ ⎛ 3 −2 2 ⎜ 2 ⎜ 0 1⎟ ⎟ ⎜ 2. A = ⎜ ⎝ 14 2⎠  B = ⎝14 1 6 8    x 1−x 1 3. A = B = 2 ex x

requested ⎞ 0 ⎟ 6⎠  4

⎞ 4 1⎟ ⎟  −5A + 3B 16⎠ 25  −6  A2 + 2AB cos
x

4. A =
14 B =
−12 −3A − 5B   1 −2 1 7 −9  5. A = 8 2 −5 0 0   −5 1 8 21 7  4A + 8B B= 12 −6 −2 −1 9     0 8 −2 3 B =  A3 − B2 6. A = −5 1 1 1 In each of Problems 7 through 16, determine which of AB and BA are defined. Carry out the products that are defined.

⎛

−4 7. A = ⎝−2 1  −2 8. A = 3  9. A = −1 ⎛

10.

11.

12.

13.

−3 ⎜ 6 A=⎜ ⎝ 18 1 ⎛ −21 ⎜ 12 A=⎜ ⎝ 1 13  −2 A= 3 ⎛ −4 A=⎝ 0 −3

⎞ ⎞ ⎛ −2 4 6 12 5 6 2 1 4⎠ −2 3⎠  B = ⎝−3 −3 1 0 0 1 6 −9 1 8    −4 6 8 B = 1 −4 −1 ⎛ ⎞ −3 ⎜ 2⎟ ⎜ ⎟  ⎟ 6 2 14 −22  B = ⎜ ⎜ 6⎟ ⎝ 0⎠ −4 ⎞ 1   2⎟ ⎟  B = −16 0 0 28 −22⎠ 0 1 1 26 6 ⎞ 4 8 −3   3 2 1 0 14⎟ ⎟  B = −9 16 5 9 14 0 16 0 −8⎠ 4 8 0    1 −3 7 2 4 B = −5 6 1 0 9 ⎞ −2 0   5 3⎠  B = 1 −3 4 1 1

7.2 Elementary Row Operations and Elementary Matrices ⎞ 3 ⎜ 0⎟   ⎟ 14. A = ⎜ ⎝−1⎠  B = 3 −2 7 4    7 −8 1 −4 B = 15. A = −4 7 1 6 ⎞ ⎛ −3 2 ⎜ 0 −2⎟  ⎟  B = −5 5 ⎜ 16. A = ⎝ 1 8⎠ 3 −3 ⎛

251

24. Let G be the graph of Figure 7.5. Determine the number of v1 − v4 walks of length 4 in G. Determine the number of v2 − v3 walks of length 2. 3 0

7



v1 2



In each of Problems 17 through 21,determine if AB is defined and if BA is defined. For those products that are defined, give the dimensions of the product matrix. 17. A is 14 × 21, B is 21 × 14. 18. A is 18 × 4, B is 18 × 4. 19. A is 6 × 2, B is 4 × 6.

v2 v5 v3

v4

FIGURE 7.5

25. Let G be the graph of Figure 7.6. Determine the number of v4 − v5 walks of length 2, the number of v2 − v3 walks of length 3, and the number of v1 − v2 and v4 − v5 walks of length 4 in G.

20. A is 1 × 3, B is 3 × 3. 21. A is 7 × 6, B is 7 × 7. 22. Find nonzero 2 × 2 matrices A, B, and C such that BA = CA but B = C. 23. Let G be the graph of Figure 7.4. Determine the number of v1 − v4 walks of length 3, the number of v2 − v3 walks of length 3, and the number of v2 − v4 walks of length 4 in G.

v5

v1 v2 v4

v3

FIGURE 7.6

26. Let A be the adjacency matrix of a graph G.

v2 v1 v3

v5 v4 FIGURE 7.4

7.2

(a) Prove that the i j− element of A2 equals the number of points of G that are neighbors of vi in G. This number is called the degree of vi . (b) Prove that the i j− element of A3 equals twice the number of triangles in G containing vi as a vertex. A triangle in G consists of three points, each a neighbor of the other two.

Elementary Row Operations and Elementary Matrices When we solve a system of linear algebraic equations by elimination of unknowns, we routinely perform three kinds of operations: interchange of equations, multiplication of an equation by a nonzero constant, and addition of a constant multiple of one equation to another equation. When we write a homogeneous system in matrix form AX = O, row k of A lists the coefficients in equation k of the system. The three operations on equations correspond, respectively, to the interchange of two rows of A, multiplication of a row A by a constant, and addition of a scalar multiple of one row of A to another row of A. We will focus on these row operations in anticipation of using them to solve the system.

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DEFINITION 7.9

Let A be an n × m matrix. The three elementary row operations that can be performed on A are: 1. Type I operation: interchange two rows of A. 2. Type II operation: Multiply a row of A by a nonzero constant. 3. Type III operation: Add a scalar multiple of one row to another row.

The rows of A are m-vectors. In a Type II operation, multiply a row by a nonzero constant by multiplying this row vector by the number. That is, multiply each element of the row by that number. Similarly, in a Type III operation, we add a scalar multiple of one row vector to another row vector.

EXAMPLE 7.8

Let ⎛

−2 ⎜ 1 A=⎜ ⎝ 0 2

1 1 1 −3

Type I operation: if we interchange rows 2 and 4 ⎛ −2 1 ⎜ 2 −3 ⎜ ⎝ 0 1 1 1

⎞ 6 2⎟ ⎟ 3⎠ 4

of A, we obtain the new matrix ⎞ 6 4⎟ ⎟ 3⎠ 2

Type II operation: multiply row 2 of A by 7 to get ⎛ ⎞ −2 1 6 ⎜ 7 7 14⎟ ⎜ ⎟ ⎝ 0 1 3⎠ 2 −3 4 Type III operation: add 2 times row 1 to row 3 of A, obtaining ⎛ ⎞ −2 1 6 ⎜ 1 1 2⎟ ⎜ ⎟ ⎝−4 3 15⎠ 2 −3 4 Elementary row operations can be performed on any matrix. When performed on an identity matrix, we obtain special matrices that will be particularly useful. We therefore give matrices formed in this way a name.

7.2 Elementary Row Operations and Elementary Matrices

DEFINITION 7.10

253

Elementary Matrix

An elementary matrix is a matrix formed by performing an elementary row operation on In .

For example,

⎛

0 ⎝1 0

1 0 0

⎞ 0 0⎠ 1

is an elementary matrix, obtained from I3 by interchanging rows 1 and 2. And ⎛ ⎞ 1 0 0 ⎝ 0 1 0⎠ −4 0 1 is the elementary matrix formed by adding −4 times row 1 of I3 to row 3. The following theorem is the reason why elementary matrices are interesting. It says that each elementary row operation on A can be performed by multiplying A on the left by an elementary matrix. THEOREM 7.5

Let A be an n × m matrix. Let B be formed from A by an elementary row operation. Let E be the elementary matrix formed by performing this elementary row operation on In . Then B = EA We leave a proof to the exercises. It is instructive to see the theorem in practice.

EXAMPLE 7.9

Let

⎛

⎞ 1 −5 4⎠  A=⎝ 9 −3 2

Suppose we form B from A by interchanging rows 2 and 3 of A. We can do this directly. But we can also form an elementary matrix by performing this operation on I3 to form ⎛ ⎞ 1 0 0 E = ⎝0 0 1⎠  0 1 0 Then

⎛

1 EA = ⎝0 0

0 0 1

⎞ ⎞ ⎛ ⎞⎛ 1 −5 1 −5 0 2⎠ = B 4⎠ = ⎝−3 1⎠ ⎝ 9 9 4 −3 2 0

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EXAMPLE 7.10

Let A=

 0 5

 −7 3 6  1 −11 3

Form C from A by multiplying row 2 by −8. Again, we can do this directly. However, if we form E from by performing this operation on I2 , then   1 0 E= 0 −8 and



  1 0 0 −7 3 6 EA = 0 −8 5 1 −11 3   0 −7 3 6 = = C −40 −8 88 −24

EXAMPLE 7.11

Let

⎛

−6 14 4 A=⎝ 4 −3 2 Form D from A by adding 6 times row 1 to row 2. ⎛ 1 0 E = ⎝6 1 0 0 then

⎞ 2 −9⎠  13 If we perform this operation on I3 to form ⎞ 0 0⎠  1

⎞ ⎞⎛ −6 14 2 0 4 −9⎠ 0⎠ ⎝ 4 1 −3 2 13 ⎞ −6 14 2 3⎠ = B = ⎝−32 88 −3 2 13 ⎛

1 EA = ⎝6 0 ⎛

0 1 0

Later we will want to perform not just one elementary row operation, but a sequence of such operations. Suppose we perform operation 1 on A to form A1 , then operation 2 on A1 to form A2 , and so on until finally we perform r on Ar−1 to get Ar . This process can be diagrammed: A → A1 → A2 → · · · → Ar−1 → Ar  1

2

3

r−1

r

7.2 Elementary Row Operations and Elementary Matrices

255

Let Ej be the elementary matrix obtained by performing operation j on In . Then A1 = E1 A A2 = E2 A1 =
E2 E1 A A3 = E3 A2 =
E3 E2 E1 A   Ar = Er Ar−1 =
Er Er−1 · · · E3 E2 E1 A This forms a matrix  = Er Er−1 · · · E2 E1 such that Ar = A

The significance of this equation is that we have produced a matrix  such that multiplying A on the left by  performs a given sequence of elementary row operations.  is formed as a product Er Er−1    E2 E1 of elementary matrices, in the correct order, with each elementary matrix performing one of the prescribed elementary row operations in the sequence (E1 performs the first operation, E2 the second, and so on until Er performs the last). We will record this result as a theorem. THEOREM 7.6

Let A be an n × m matrix. If B is produced from A by any finite sequence of elementary row operations, then there is an n × n matrix  such that B = A The proof of the theorem is contained in the line of reasoning outlined just prior to its statement.

EXAMPLE 7.12

Let

⎛

⎞ 2 1 0 A = ⎝ 0 1 2⎠  −1 3 2 We will form a new matrix B from A by performing, in order, 1 : interchange rows 1 and 2 of A to form A1 . 2 : multiply row 3 of A1 by 2 to form A2 . 3 : add two times row 1 to row 3 of A2 to get A3 = B If we perform this sequence in order, starting with A, we get ⎛ ⎞ ⎛ 0 1 2 0 1 A → A1 = ⎝ 2 1 0⎠ → A2 = ⎝ 2 1 1 2 −1 3 2 −2 6 ⎛ ⎞ 0 1 2 → A3 = ⎝ 2 1 0⎠ = B 3 −2 8 8

the following operations:

⎞ 2 0⎠ 4

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Matrices and Systems of Linear Equations

To produce  such that B = A, perform this sequence of operations in turn, beginning with I3 : ⎛ ⎞ ⎛ ⎞ 0 1 0 0 1 0 I3 → ⎝ 1 0 0 ⎠ → ⎝ 1 0 0 ⎠ 1 2 0 0 1 0 0 2 ⎛ ⎞ 0 1 0 → ⎝1 0 0⎠ = . 3 0 2 2 Now check that

⎛

0 A = ⎝1 0 ⎛

0 =⎝ 2 −2

1 0 2

⎞⎛ ⎞ 0 2 1 0 0⎠ ⎝ 0 1 2⎠ 2 −1 3 2 ⎞ 1 2 1 0⎠ = B 8 8

It is also easy to check that  = E3 E2 E1 , where Ej is the elementary matrix obtained by performing operation j on I3 .

EXAMPLE 7.13

Let

⎛

6 A = ⎝9 0

−1 1 3 7 2 1

⎞ 4 −7⎠  5

We want to perform, in succession and in the given order, the following operations: 1 : add
−3
row 2) to row 3, 2 : add 2(row 1) to row 2, 3 : interchange rows 1 and 3, 4 : multiply row 2 by −4. Suppose the end result of these operations is the matrix B. We will produce a 3 × 3 matrix  such that B = A. Perform the sequence of operations, starting with I3 : ⎛ ⎞ ⎛ ⎞ 1 0 0 1 0 0 1 0 ⎠ → ⎝2 1 0⎠ I3 → ⎝0 1 2 0 −3 1 0 −3 1 ⎛ ⎞ ⎛ ⎞ 0 −3 1 0 −3 1 1 0⎠ → ⎝−8 −4 0⎠ =  → ⎝2 3 4 1 0 0 1 0 0 Then

⎛

⎞⎛ ⎞ 0 −3 1 6 −1 1 4 3 7 −7⎠ A = ⎝−8 −4 0⎠ ⎝9 1 0 0 0 2 1 5 ⎛ ⎞ −27 −7 −20 26 = ⎝−84 −4 −36 −4⎠ = B 6 −1 1 4

7.2 Elementary Row Operations and Elementary Matrices

257

It is straightforward to check that  = E4 E3 E2 E1 , where Ej is the elementary matrix obtained from I3 by applying operation j . If the operations j are performed in succession, starting with A, then B results.

DEFINITION 7.11

Row Equivalence

Two matrices are row equivalent if and only if one can be obtained from the other by a sequence of elementary row operaitons.

In each of the last two examples, B is row equivalent to A. The relationship of row equivalence has the following properties: THEOREM 7.7

1. Every matrix is row equivalent to itself. (This is the reflexive property). 2. If A is row equivalent to B, then B is row equivalent to A. (This is the symmetry property). 3. If A is row equivalent to B, and B to C, then A is row equivalent to C. (This is transitivity). It is sometimes of interest to undo the effect of an elementary row operation. This can always be done by the same kind of elementary row operation. Consider each kind of operation in turn. If we interchange rows i and j of A to form B, then interchanging rows i and j of B yields A again. Thus a Type I operation can reverse a Type I operation. If we form C from A by multiplying row i by nonzero , then multiplying row i of C by 1/ brings us back to A. A Type II operation can reverse a Type II operation. Finally, suppose we form D from A by adding  (row i) to row j. Then ⎛ ⎞ a11 a12 · · · a1m ⎜· · · · · · · · · · · · ⎟ ⎜ ⎟ ⎜ ai1 ai2 · · · aim ⎟ ⎜ ⎟ ⎟ A=⎜ ⎜· · · · · · · · · · · · ⎟ ⎜ aj1 aj2 · · · ajm ⎟ ⎜ ⎟ ⎝· · · · · · · · · · · · ⎠ an1 an2 · · · anm and

⎛

a11 ⎜ ··· ⎜ ⎜ ai1 ⎜ D=⎜ ⎜ ··· ⎜ai1 + aj1 ⎜ ⎝ ··· an1

a12 ··· ai2 ··· ai2 + aj2 ··· an2

⎞ ··· a1m ⎟ ··· ··· ⎟ ⎟ ··· aim ⎟ ⎟ ··· ··· ⎟ · · · aim + ajm ⎟ ⎟ ⎠ ··· ··· ··· anm

Now we can get from D back to A by adding −(row i) to row j of D. Thus a Type III operation can be used to reverse a Type III operation. This ability to reverse the effects of elementary row operations will be useful later, and we will record it as a theorem.

CHAPTER 7

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Matrices and Systems of Linear Equations

THEOREM 7.8

Let E1 be an elementary matrix that performs an elementary row operation on a matrix A. Then there is an elementary matrix E2 such that E2
E1 A = A. In fact, E2 E1 = In .

SECTION 7.2

PROBLEMS

In each of Problems 1 through 8, perform the row operation, or sequence of row operations, directly on A, and then find a matrix  such that the final result is A. ⎞ ⎛ −2 1 4 2 √ 1 16 3⎠; multiply row 2 by 3 1. A = ⎝ 0 1 −2 4 8 ⎞ ⎛ 3 −6 ⎜1 1⎟ ⎟; add 6 times row 2 to row 3. ⎜ 2. A = ⎝ 8 −2⎠ 0 5 ⎞ ⎛ −2 14 6 √ 1 −3⎠; add 13 times row 3 to row 1, 3. A = ⎝ 8 2 9 5 then interchange rows 2 and 1, then multiply row 1 by 5. ⎞ ⎛ −4 6 −3 4. A = ⎝ 12 4 −4⎠; interchange rows 2 and 3, then 1 3 0 add negative row 1 to row 2.   √ −3 15 ; add 3 times row 2 to row 1, then 5. A = 2 8 multiply row 2 by 15, then interchange rows 1 and 2. ⎞ ⎛ 3 −4 5 9 1 3 −6⎠; add row 1 to row 3, then 6. A = ⎝2 1 13 2 6

7.3

√ add 3 times row 1 to row 2, then multiply row 3 by row 4, then add row 2 to row 3. ⎞ ⎛ −1 0 3 0 2 9⎠; multiply row 3 by 4, then 7. A = ⎝ 1 3 −9 7 −5 7 add 14 times row 1 to row 2, then interchange rows 3 and 2. ⎞ ⎛ 0 −9 14 5 2⎠; interchange rows 2 and 3, then 8. A = ⎝1 9 15 0 add 3 times row 2 to row 3, then interchange rows 1 and 3, then multiply row 3 by 4. In Problems 9, 10, and 11, A is an n × m matrix. 9. Let B be formed from A by interchanging rows s and t. Let E be formed from In by interchanging rows s and t. Prove that B = EA. 10. Let B be formed from A by multiplying row s by , and let E be formed from In by multiplying row s by . Prove that B = EA. 11. Let B be formed from A by adding  times row s to row t. Let E be formed from In by adding  times row s to row t. Prove that B = EA.

The Row Echelon Form of a Matrix Sometimes a matrix has a special form that makes it convenient to work with in solving certain problems. For solving systems of linear algebraic equations, we want the reduced row echelon form, or reduced form, of a matrix. Let A be an n × m matrix. A zero row of A is a row having each element equal to zero. If at least one element of a row is nonzero, that row is a nonzero row. The leading entry of a nonzero row is its first nonzero element, reading from left to right. For example, if ⎛ ⎞ 0 2 7 ⎜1 −2 0⎟ ⎟ A=⎜ ⎝0 0 0⎠ 0 0 9

7.3 The Row Echelon Form of a Matrix

259

then row three is a zero row and rows one, two, and four are nonzero rows. The leading entry of row 1 is 2, the leading entry of row 2 is 1, and the leading entry of row 4 is 9. We do not speak of a leading entry of a zero row. We can now define a reduced row echelon matrix.

DEFINITION 7.12

Reduced Row Echelon Matrix

A matrix is in reduced row echelon form if it satisfies the following conditions: 1. The leading entry of any nonzero row is 1. 2. If any row has its leading entry in column j, then all other elements of column j are zero. 3. If row i is a nonzero row and row k is a zero row, then i < k. 4. If the leading entry of row r1 is in column c1 , and the leading entry of row r2 is in column c2 , and if r1 < r2 , then c1 < c2 . A matrix in reduced row echelon form is said to be in reduced form, or to be a reduced matrix.

A reduced matrix has a very special structure. By condition (1), if we move from left to right along a nonzero row, the first nonzero number we see is 1. Condition (2) means that, if we stand at the leading entry 1 of any row, and look straight up or down, we see only zeros in the rest of this column. A reduced matrix need not have any zero rows. But if there is a zero row, it must be below any nonzero row. That is, all the zero rows are at the bottom of the matrix. Condition (4) means that the leading entries move downward to the right as we look at the matrix.

EXAMPLE 7.14

The following four matrices are all reduced:

⎞ 0 1 3 0 1 −4 1 0 ⎝  0 0 0 1⎠  0 0 0 1 0 0 0 0 ⎞ ⎛ 2 0 0 1 0 0 3 ⎜0 1 0 −2 0 1 0⎟ ⎟  and ⎜ ⎝0 0 1 0 0 0⎠ 0 0 0 0 0 0 0 0 

⎛

0 ⎜0 ⎜ ⎝0 0

1 0 0 0



⎛

⎞ 1 4⎟ ⎟ 1⎠ 0

EXAMPLE 7.15

To see one context in which reduced matrices are interesting, consider the last matrix of the preceding example and suppose it is the matrix of coefficients of a system of homogeneous linear equations. This system is AX = O, and the equations are: x1 + 3x4 + x5 = 0 x2 − 2x4 + 4x5 = 0 x3 + x5 = 0

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The fourth row represents the equation 0x1 + 0x2 + 0x3 + 0x4 = 0, which we do not write out (it is satisfied by any numbers x1 through x4 , and so provides no information). Because the matrix of coefficients is in reduced form, this system is particularly easy to solve. From the third equation, x3 = −x5  From the second equation, x2 = 2x4 − 4x5  And from the first equation, x1 = −3x4 − x5  We can therefore choose x4 = , any number, and x5 = , any number, and obtain a solution by choosing the other unknowns as x1 = −3 − 

x2 = 2 − 4

x3 = −

The form of the reduced matrix is selected just so that as a matrix of coefficients of a system of linear equations, the solution of these equations can be read by inspection.

EXAMPLE 7.16

The matrix

⎛

0 ⎜0 A=⎜ ⎝0 0

1 0 0 0

5 1 0 0

0 0 1 0

⎞ 0 0⎟ ⎟ 0⎠ 1

is not reduced. The leading entry of row 2 is 1, as it must be, but there is a nonzero element in the column containing this leading entry. However, A is row equivalent to a reduced matrix. If we add −5 (row 2) to row 1, we obtain ⎛ ⎞ 0 1 0 0 0 ⎜0 0 1 0 0⎟ ⎟ B=⎜ ⎝0 0 0 1 0⎠  0 0 0 0 1 and this is a reduced matrix.

EXAMPLE 7.17

The matrix

⎛

2 C = ⎝0 1

0 1 0

⎞ 0 0⎠ 1

is not reduced. The leading entry of the first row is not 1, and the first column, containing this leading entry of row 1, has another nonzero element. In addition, the leading entry of row 3 is to the left of the leading entry of row 2, and this violates condition (4). However, C is row

7.3 The Row Echelon Form of a Matrix equivalent to a reduced matrix. First form D by ⎛ 1 D = ⎝0 1

261

multiplying row 1 by 21 : ⎞ 0 0 1 0⎠  0 1

Now form F from D by adding − (row 1) to row 3: ⎞ ⎛ 1 0 0 F = ⎝0 1 0⎠  0 0 1 Then F is a reduced matrix that is row equivalent to C, since it was formed by a sequence of elementary row operations, starting with C. In the last two examples we had matrices that were not in reduced form, but could in both cases proceed to a reduced matrix by elementary row operations. We claim that this is always possible (although in general more operations may be needed than in these two examples). THEOREM 7.9

Every matrix is row equivalent to a reduced matrix. The proof consists of exhibiting a sequence of elementary row operations that will produce a reduced matrix. Let A be any matrix. If A is a zero matrix, we are done. Thus suppose that A has at least one nonzero row. Reading from left to right across the matrix, find the first column having a nonzero element. Suppose this is in column c1 . Reading from top to bottom in this column, suppose  is the top nonzero element. Say  is in row r1  Multiply this row by 1/ to obtain a matrix B in which column c1 has its top nonzero element equal to 1, and this is in row r1 . If any row below r1 in B has a nonzero element  in column c1 , add − times row r1 to this row. In this way we obtain a matrix C that is row equivalent to A, having 1 in the r1  c1 position, and all other elements of column c1 equal to zero. Now interchange, if necessary, rows 1 and r1 of C to obtain a matrix D having leading entry 1 in row 1 and column c1 , and all other elements of this column equal to zero. Further, by choice of c1 , any column of D to the left of column c1 has all zero elements (if there is such a column). D is row equivalent to A. If D is reduced, we are done. If not, repeat this procedure, but now look for the first column, say column c2 , to the right of column c1 having a nonzero element below row 1. Let ! be the top nonzero element of this column lying below row 1. Say this element occurs in row r2 . Multiply row r2 by 1/! to obtain a new matrix E having 1 in the r2  c2 position. If this column has a nonzero element  above or below row r2 , add − (row r2 ) to this row. In this way we obtain a matrix F that is row equivalent to A and has leading entry 1 in row r2 and all other elements of column c2 equal to zero. Finally, form G from F by interchanging rows r2 and 2, if necessary. If G is reduced we are done. If not, locate the first column to the right of column c2 having a nonzero element and repeat the procedure done to form the first two rows of G. Since A has only finitely many columns, eventually this process terminates in a reduced matrix R. Since R was obtained from A by elementary row operations, R is row equivalent to A and the proof is complete. Proof

The process of obtaining a reduced matrix row equivalent to a given matrix A is referred to as reducing A. It is possible to reduce a matrix in many different ways (that is, by different sequences of elementary row operations). We claim that this does not matter and that for a given A any reduction process will result in the same reduced matrix.

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THEOREM 7.10

Let A be a matrix. Then there is exactly one reduced matrix AR that is row equivalent to A. We leave a proof of this result to the student. In view of this theorem, we can speak of the reduced form of a given matrix A. We will denote this matrix AR .

EXAMPLE 7.18

Let

⎛

−2 A=⎝ 0 2

1 1 0

⎞ 3 1⎠  1

We want to find AR . Column 1 has a nonzero element in row 1. Begin with: ⎛ ⎞ −2 1 3 A = ⎝ 0 1 1⎠  2 0 1 Begin with the operations:

⎛

1

1 ⎜ → ⎝0 2 2 ⎛ 1 ⎜ → add
−2(row 1) to row 3 → ⎝0 0 multiply row 1 by −

− 21 1 0 − 21 1 1

− 23

⎞

⎟ 1⎠ 1 ⎞ 3

−2

⎟ 1⎠  4

In the last matrix, column 2 has a nonzero element below row 1, the highest being 1 in the 2 2 position. Since we want a 1 here, we do not have to multiply this row by anything. However, we want zeros above and below this 1, in the 1 2 and 3 2 positions. Thus add 1/2 times row 2 to row 1, and − row 2 to row 3 in the last matrix to obtain ⎛ ⎞ 1 0 −1 ⎝0 1 1⎠  0 0 3 In this matrix column 3 has a nonzero element below row 2, in the 3 3 location. Multiply row 3 by 1/3 to obtain ⎛ ⎞ 1 0 −1 ⎝0 1 1⎠  0 0 1 Finally, we want zeros above the 3 3 position in column 3. Add row 3 to row 1 and − row 3 to row 2 to get ⎛ ⎞ 1 0 0 AR = ⎝0 1 0⎠  0 0 1 This is AR because it is a reduced matrix and it is row equivalent to A.

7.3 The Row Echelon Form of a Matrix

263

To illustrate the last theorem, we will use a different sequence of elementary row operations to reduce A, arriving at the same final result. Proceed: ⎛ 0 A → (add row 3 to row 1) → ⎝0 2 ⎛ 0 add
−1(row 2) to row 1 → ⎝0 2 ⎛ 0 1 (row 1) → ⎝0 3 2

⎞ 4 1⎠ 1 ⎞ 3 1⎠ 1 ⎞ 1 1⎠ 1 ⎞ 1 0⎠ 0 ⎞ 0 0⎠ 1 ⎞ 0 0 ⎠ = AR  1

1 1 0 0 1 0 0 1 0

⎛

0 add
−1(row 1) to rows 2 and 3 → ⎝0 2 ⎛ 2 interchange rows 1 and 3 → ⎝0 0 ⎛ 1 1 (row 1) → ⎝0 2 0

0 1 0 0 1 0 0 1 0

EXAMPLE 7.19

Let ⎛

0 ⎜0 B=⎜ ⎝0 0

0 0 1 4

0 2 0 3

0 0 1 4

⎞ 0 0⎟ ⎟ 1⎠ 0

Reduce B as follows: ⎛

0 ⎜0 add − 4(row 3) to row 4 → ⎜ ⎝0 0 ⎛ 0 ⎜0 interchange rows 3 and 1 → ⎜ ⎝0 0 ⎛ 0 ⎜0 1 (row 2) → ⎜ ⎝0 2 0

0 0 1 0

0 2 0 3

0 0 1 0

1 0 0 0

0 2 0 3

1 0 0 0

1 0 0 0

0 1 0 3

1 0 0 0

⎞ 0 0⎟ ⎟ 1⎠ −4 ⎞ 1 0⎟ ⎟ 0⎠ −4 ⎞ 1 0⎟ ⎟ 0⎠ −4

264

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Matrices and Systems of Linear Equations ⎛

0 ⎜0 add
−3
row 2) to row 4 → ⎜ ⎝0 0 ⎛ 0 ⎜0 1 − (row 4) → ⎜ ⎝0 4 0 ⎛ 0 ⎜0 add
−1(row 4) to row 1 → ⎜ ⎝0 0 ⎛ 0 ⎜0 interchange row 3 and 4 → ⎜ ⎝0 0

1 0 0 0

0 1 0 0

1 0 0 0

1 0 0 0

0 1 0 0

1 0 0 0

1 0 0 0

0 1 0 0

1 0 0 0

1 0 0 0

0 1 0 0

1 0 0 0

⎞ 1 0⎟ ⎟ 0⎠ −4 ⎞ 1 0⎟ ⎟ 0⎠ 1 ⎞ 0 0⎟ ⎟ 0⎠ 1 ⎞ 0 0⎟ ⎟ 1⎠ 0

This is a reduced matrix, hence it is the reduced matrix BR of B. In view of Theorem 7.6 of the preceding section, we immediately have the following.

THEOREM 7.11

Let A be an n × m matrix. Then there is an n × n matrix  such that A = AR . There is a convenient notational device that enables us to find both  and AR together. We know what  is. If A is n = n × m, then  is an n × n matrix formed by starting with In and carrying out, in order, the same sequence of elementary row operations used  to reduce A. A simple way to form  while reducing A is to form an n ×
n + m matrix

I A by placing n  In alongside A on its left. The first n columns of this matrix In A are just In , and the last m columns are A. Now reduce A by elementary row operations, performing the same operations on the first n columns (In ) as well. When A is reduced, the resulting n ×
n + m matrix will  have the form AR , and we read  as the first n columns.

EXAMPLE 7.20

Let  A=

 −3 1 0  4 −2 1

We want to find a 2 × 2 matrix  such that A = AR . Since A is 2 × 3, form the matrix ⎛  1

I2 A = ⎝ 0

0 1

   

−3 4

⎞ 1 0⎠  −2 1

7.3 The Row Echelon Form of a Matrix

265

Now reduce the last three columns, performing the same operations on the first two. The column of dashes is just a bookkeeping device to separate A from I2 . Proceed ⎛ ⎞  1 1  1 − 0  1 − 0 3 ⎠

I2 A → − (row 1) → ⎝ 3  3 0 1  4 −2 1 ⎛ 1 −1 → (row 2) → ⎝ 3 4 0 ⎛ −1 → (row 2 − row 1) → ⎝ 3

1 4

→
−6(row 2) →

   

0

1 3

1 4

⎛ 1 ⎝− 3

− 23

⎛ 1 −1 → (row 2) + (row 1) → ⎝ 3 −2

   

− 21 − 23

1 4

1

− 13

⎞ 0⎠

0

− 16

1 4

   

0

−2

⎞ 0⎠

  1 − 13   1 − 21

0

1

− 13

0

1

⎞ 0⎠ − 23 ⎞

1

0

− 21 ⎠

0

1

− 23



The last three columns are in reduced form, so they form AR . The first two columns form :   −1 − 21 =  −2 − 23 As a check on this, form the product   −1 − 21 −3 A = 4 −2 − 23

SECTION 7.3

  1 1 0 = −2 1 0

0

− 21

1

− 23

PROBLEMS

In each of the following, find the reduced form of A and produce a matrix  such that A = AR . ⎞ ⎛ 1 −1 3 1 2⎠ 1. A = ⎝0 0 0 0   3 1 1 4 2. A = 0 1 0 0 ⎞ ⎛ −1 4 1 1 ⎜ 0 0 0 0⎟ ⎟ 3. A = ⎜ ⎝ 0 0 0 0⎠ 0 0 0 1   1 0 1 1 −1 4. A = 0 1 0 0 2

⎛

6 ⎜0 5. A = ⎜ ⎝1 0 

2 1

6. A =

⎛

−1 7. A = ⎝ 2 7  8. A =

−3 0

⎞ 1 0⎟ ⎟ 3⎠ 1 2 1



4 3 1

⎞ 6 −5⎠ 1

4 0

4 0



 = AR 

CHAPTER 7

266

 9. A =

−1 1

⎛

8 10. A = ⎝0 4

7.4

2 0 2 1 0

3 0 1 1 0

Matrices and Systems of Linear Equations 1 0



⎛

4 1 11. A = ⎝2 2 0 1 ⎛ ⎞ 6 ⎜−3⎟ ⎟ 12. A = ⎜ ⎝ 1⎠ 1

⎞ 0 3⎠ −3

⎞ −7 0⎠ 0

The Row and Column Spaces of a Matrix and Rank of a Matrix In this section we will develop three numbers associated with matrices that play a significant role in the solution of systems of linear equations. Suppose A is an n × m matrix with real number elements. Each row of A has m elements and can be thought of as a vector in Rm . There are n such vectors. The set of all linear combinations of these row vectors is a subspace of Rm called the row space of A. This space is spanned by the row vectors. If these row vectors are linearly independent, they form a basis for this row space and this space has dimension n. If they are not linearly independent, then some subset of them forms a basis for the row space, and this space has dimension < n. If we look down instead of across, we can think of each column of A as a vector in Rn . We often write these vectors as columns simply to keep in mind their origin, although they can be written in standard vector notation. The set of all linear combinations of these columns forms a subspace of Rn . This is the column space of A. If these columns are linearly independent, they form a basis for this column space, which then has dimension m; otherwise, this dimension is less than m.

EXAMPLE 7.21

Let

⎛

⎞ −2 6 1 ⎜ 2 2 −4⎟ ⎜ ⎟ ⎜ B = ⎜ 10 −8 12⎟ ⎟ ⎝ 3 1 −2⎠ 5 −5 7

The row space is the subspace of R3 spanned by the row vectors of B. This row space consists of all vectors 
−2 6 1 + 
2 2 −4 + !
10 −8 12 + 
3 1 −2 + 
5 −5 7 The first three row vectors are linearly independent. The last two are linear combinations of the first three. Specifically,
3 1 −2 =

4 181 13
−2 6 1 +
2 2 −4 +
10 −8 12 101 202 101

and
5 −5 7 = −

39 53 7
−2 6 1 −
2 2 −4 +
10 −8 12 101 202 101

The first three row vectors form a basis for the row space, which therefore has dimension 3. The row space of B is all of R3 .

7.4 The Row and Column Spaces of a Matrix and Rank of a Matrix

267

The column space of B is the subspace of R5 consisting of all vectors ⎛

⎞ ⎛ ⎞ ⎛ ⎞ −2 6 1 ⎜ 2⎟ ⎜ 2⎟ ⎜−4⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎜ 10⎟ +  ⎜−8⎟ + ! ⎜ 12⎟  ⎝ 3⎠ ⎝ 1⎠ ⎝−2⎠ 5 −5 7 These three column vectors are linearly independent in R5 . Neither is a linear combination of the other two, or, equivalently, the only way this linear combination can be the zero vector is for  =  = ! = 0. Therefore the column space of B has dimension 3 and is a subspace of the dimension 5 space R5 . In this example the row space of matrix had the same dimension as the column space, even though the row vectors were in Rm and the column vectors in Rn , with n = m. This is not a coincidence.

THEOREM 7.12

For any matrix A having real numbers as elements, the row and column spaces have the same dimension. Proof

Suppose A is n × m: ⎛

a11 a21  

⎜ ⎜ ⎜ ⎜ ⎜ A=⎜ ⎜ ar1 ⎜ar+11 ⎜ ⎜  ⎝  an1

a12 a22  

··· ···  

a1r a2r  

a1r+1 a2r+1  

··· ···  

ar2 ar+12  

··· ···  

arr ar+1r  

arr+1 ar+1r+1  

··· ···  

an2

···

anr

anr +1

···

a1m a2m  

⎞

⎟ ⎟ ⎟ ⎟ ⎟ arm ⎟ ⎟ ar+1m ⎟ ⎟  ⎟  ⎠ anm

Denote the row vectors R1      Rn , so Ri =
ai1  ai2  · · ·  aim  in Rm  Now suppose that the dimension of the row space of A is r. Then exactly r of these row vectors are linearly independent. As a notational convenience, suppose the first r rows R1      Rr are linearly independent. Then each of Rr+1      Rn is a linear combination of these r vectors. Write Rr+1 = r+11 R1 + · · · + r+1r Rr  Rr+2 = r+21 R1 + · · · + r+2r Rr    Rn = n1 R1 + · · · + nr Rr 

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Now observe that column j of A can be written ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ a1j 1 0 0 ⎜ 0 ⎟ ⎜ 1 ⎟ ⎜ 0 ⎟ ⎜ a2j ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ arj ⎟ = a1j ⎜ 0 ⎟ + a2j ⎜ 0 ⎟ + · · · + arj ⎜ 1 ⎟  ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ar+1j ⎟ ⎜r+11 ⎟ ⎜r+12 ⎟ ⎜r+1r ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎜  ⎟ ⎝  ⎠ ⎝  ⎠ ⎝  ⎠ ⎝  ⎠ anj n1 n2 nr This means that each column vector of A is a linear combination of the r n-vectors on the right side of this equation. These r vectors therefore span the column space of A. If these vectors are linearly independent, then the dimension of the column space is r. If not, then remove from this list of vectors any that are linear combinations of the others and thus determine a basis for the column space having fewer than r vectors. In any event, Dimension of the column space of A ≤ dimension of the row space of A. By essentially repeating this argument, with row and column vectors interchanged, we obtain Dimension of the row space of A ≤ dimension of the column space of A and these two inequalities together prove the theorem.

It is interesting to ask what effect elementary row operations have on the row space of a matrix. The answer is—none! We will need this fact shortly. THEOREM 7.13

Let A be an n × m matrix, and let B be formed from A by an elementary row operation. Then the row space of A and the row space of B are the same. If B is obtained by a Type I operation, we simply interchange two rows. Then A and B still have the same row vectors, just listed in a different order, so these row vectors span the same row space. Suppose B is obtained by a Type II operation, multiplying row i by a nonzero constant c. Linear combinations of the rows of A have the form

Proof

1 R1 + · · · + i Ri + · · · + n Rn while linear combinations of the rows of B are 1 R1 + · · · + ci Ri + · · · + n Rn Since i can be any number, so can ci , so these linear combinations yield the same vectors when the coefficients are chosen arbitrarily. Thus the row space of A and B are again the same. Finally, suppose B is obtained from A by adding c (row i) to row j. The column vectors of B are now R1     Rj−1  cRi + Rj  Rj+1      Rn 

7.4 The Row and Column Spaces of a Matrix and Rank of a Matrix

269

But we can write an arbitrary linear combination of these rows of B as   1 R1 + · · · + j−1 Rj−1 + j cRi + Rj + j+1 Rj+1 + · · · + n Rn  and this is 1 R1 + · · · +
i + cj Ri + · · · + j Rj + · · · + n Rn  which is again just a linear combination of the row vectors of A. Thus again the row spaces of A and B are the same, and the theorem is proved.

COROLLARY 7.1

For any matrix A, the row spaces of A and AR are the same. This follows immediately from Theorem 6.13. Each time we perform an elementary row operation on a matrix, we leave the row space unchanged. Since we obtain AR from A by elementary row operations, then A and AR must have the same row spaces. The dimensions of the row and column spaces will be important when we consider solutions of systems of linear equations. There is another number that will play a significant role in this, the rank of a matrix.

DEFINITION 7.13

Rank

The rank of a matrix A is the number of nonzero rows in AR .

We denote the rank of A as rank
A. If B is a reduced matrix, then B = BR , so the rank of B is just the number of nonzero rows of B itself. Further, for any matrix A rank
A = number of nonzero rows of AR = rank
AR  We claim that the rank of a matrix is equal to the dimension of its row space (or column space). First we will show this for reduced matrices.

LEMMA 7.1

Let B be a reduced matrix. Then the rank of B equals the dimension of the row space of B. Let R1      Rr be the nonzero row vectors of B. The row space consists of all linear combinations

Proof

c1 R1 + · · · + cr Rr  If nonzero row j has its leading entry in column k, then the kth component of Rj is 1. Because B is reduced, all the other elements of column k are zero, hence each other Ri has kth component zero. By Lemma 5.3, R1      Rr are linearly independent. Therefore these vectors form a basis for the row space of B, and the dimension of this space is r. But r = number of nonzero rows of B = number of nonzero rows of BR = rank
B.

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EXAMPLE 7.22

Let

⎛

0 ⎜0 ⎜ B=⎝ 0 0

1 0 0 0

0 1 0 0

0 0 1 0

3 0 −2 1 2 0 0 0

⎞ 6 5⎟ ⎟ −4⎠ 0

Then B is in reduced form, so B = BR . The rank of B is its number of nonzero rows, which is 3. Further, the nonzero row vectors are
0 1 0 0 3 0 6
0 0 1 0 −2 1 5
0 0 0 1 2 0 −4 and these are linearly independent. Indeed, if a linear combination of these vectors yielded the zero vector, we would have 
0 1 0 0 3 0 6 + 
0 0 1 0 −2 1 5 + !
0 0 0 1 2 0 −4 =
0 0 0 0 0 0 0 But then
0   ! 3 − 2 + 2!  6 + 5 − 4! =
0 0 0 0 0 0 0 and from the second, third, and fourth components we read that  =  = ! = 0. By Theorem 6.13(2), these three row vectors are linearly independent and form a basis for the row space, which therefore has dimension 3. Using this as a stepping stone, we can prove the result for arbitrary matrices.

THEOREM 7.14

For any matrix A, the rank of A equals the dimension of the row space of A. Proof

From the lemma, we know that rank
A = rank
AR  = dimension of the row space of AR = dimension of the row space of A,

since A and AR have the same row space. Of course, we can also assert that rank
A = dimension of the column space of A. If A is n × m, then so is AR . Now AR cannot have more than n nonzero rows (because it has only n rows). This means that rank
A ≤ number of rows of A. There is a special circumstance in which the rank of a square matrix actually equals its number of rows.

7.4 The Row and Column Spaces of a Matrix and Rank of a Matrix

271

THEOREM 7.15

Let A be an n × n matrix. Then rank
A = n if and only if AR = In . If AR = In , then the number of nonzero rows in AR is n, since In has no zero rows. Hence in this case rank
A = n. Conversely, suppose that rank
A = n. Then AR has n nonzero rows, hence no zero rows. By definition of a reduced matrix, each row of AR has leading entry 1. Since each row, being a nonzero row, has a leading entry, then the i i elements of AR are all equal to 1. But it is also required that, if column j contains a leading entry, then all other elements of that column are zero. Thus AR must have each i j element equal to zero if i = j, so AR = In .

Proof

EXAMPLE 7.23

Let

⎛

⎞ 1 −1 4 2 1 3 2⎠  A = ⎝0 3 −2 15 8

We find that

⎛

1 AR = ⎝0 0

0 1 0

7 3 0

⎞ 0 2⎠  0

Therefore rank
A = 2. This is also the dimension of the row space of A and of the column space of A.

In the next section we will use the reduced form of a matrix to solve homogeneous systems of linear algebraic equations.

SECTION 7.4

PROBLEMS

In each of Problems 1 through 14, (a) find the reduced form of the matrix, and from this the rank, (b) find a basis for the row space of the matrix, and the dimension of this space, and (c) find a basis for the column space and the dimension of this space. 

−4 2 ⎛ 1 2. ⎝0 2 ⎛ −3 3. ⎝ 2 4 1.

1 2 −1 1 −1



3 0

⎛

4.

5. 6.

⎞

4 3⎠ 11 ⎞ 1 2⎠ −3

7.

8.

6 ⎝12 1  8 1  1 0 ⎛ 2 ⎜1 ⎜ ⎝0 4 ⎛ 0 ⎝0 0

0 0 −1 −4 −1

0 0 0

3 0 2 −1 0 0 −1 0 0

0 1

3 1 

2 0

1 2 0 

⎞ 1 3⎟ ⎟ 1⎠ 7 ⎞ 0 −1⎠ 2

⎞ 1 2⎠ 0

272 ⎛

0 9. ⎝6 2 ⎛ 1 ⎜2 10. ⎜ ⎝1 3 ⎛ −3 11. ⎝ 1 0

7.5

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Matrices and Systems of Linear Equations

⎞ 3 0⎠ 2 ⎞ 0 0 0 0⎟ ⎟ 0 −1⎠ 0 0 ⎞ 2 2 0 5⎠ 0 2 4 1 2

⎛

−4 −2 1 4 −4 12. ⎝ 0 1 0 0 ⎞ ⎛ −2 5 7 13. ⎝ 0 1 −3⎠ −4 11 11  −3 2 1 14. 6 −4 −2

⎞ 6 2⎠ 0

1 −2

0 0



15. Show that for any matrix A, rank
A = rank
At .

Solution of Homogeneous Systems of Linear Equations We will apply the matrix machinery we have developed to the solution of systems of n linear homogeneous equations in m unknowns: a11 x1 + a12 x2 + · · · + a1m xm = 0 a2 x1 + a22 x2 + · · · + a2m xm = 0   an1 x1 + an2 x2 + · · · + anm xm = 0 This term homogeneous applies here because the right side of each equation is zero. As a prelude to a matrix approach to solving this system, consider the simple system x1 − 3x2 + 2x3 = 0 −2x1 + x2 − 3x3 = 0 We can solve this easily by “eliminating unknowns.” Add 2(equation 1) to equation 2 to get −5x2 + x3 = 0 hence 1 x 2 = x3  5 Now put this into the first equation of the system to get 3 x1 − x3 + 2x3 = 0 5 or 7 x1 + x3 = 0 5 Then 7 x 1 = − x3  5 We now have the solution: 7 1 x1 = −  x2 =  x3 =  5 5

7.5 Solution of Homogeneous Systems of Linear Equations

273

in which  can be any number. For this system, two of the unknowns can be written as constant multiples of the third, which can be assigned any value. The system therefore has infinitely many solutions. For this simple system we do not need matrices. However, it is instructive to see how matrices could be used here. First, write this system in matrix form as AX = O, where ⎛ ⎞   x1 1 −3 2 A= and X = ⎝x2 ⎠  −2 1 −3 x3 Now reduce A. We find that

 AR =

1

0

0

1

7 5 − 15

 

The system AR X = 0 is just 7 x1 + x3 = 0 5 1 x2 − x3 = 0 5 This reduced system has the advantage of simplicity—we can solve it on sight, obtaining the same solutions that we got for the original system. This is not a coincidence. AR is formed from A by elementary row operations. Since each row of A contains the coefficients of an equation of the system, these row operations correspond in the system to interchanging equations, multiplying an equation by a nonzero constant, and adding a constant multiple of one equation to another equation of the system. This is why these elementary row operations were selected. But these operations always result in new systems having the same solutions as the original system (a proof of this will be given shortly). The reduced system AR X = O therefore has the same solutions as AX = O. But AR is defined in just such a way that we can just read the solutions, giving some unknowns in terms of others, as we saw in this simple case. We will look at two more examples and then say more about the method in general.

EXAMPLE 7.24

Solve the system x1 − 3x2 + x3 − 7x4 + 4x5 = 0 x1 + 2x2 − 3x3 = 0 x2 − 4x3 + x5 = 0 This is the system AX = O, with

⎛

1 A = ⎝1 0 We find that

⎛

1

⎜ A R = ⎝0 0

−3 2 1

1 −3 −4

−7 0 0

0

0

− 35 16

1

0

0

1

28 16 7 16

⎞ 4 0⎠  1 13 ⎞ 16 ⎟ − 20  16 ⎠ 9 − 16

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The systems AX = O and AR X = O have the same solutions. But the equations of the reduced system AR X = O are 35 13 x + x =0 16 4 16 5 28 20 x2 + x4 − x5 = 0 16 16 7 9 x3 + x4 − x5 = 0 16 16 x1 −

From these we immediately read the solution. We can let x4 =  and x5 =  (any numbers), and then x1 =

13 20 9 35 28 7  −  x2 = −  +  x3 = −  +  16 16 16 16 16 16

Not only did we essentially have the solution once we obtained AR , but we also knew the number of arbitrary constants that appear in the solution. In the last example this number was 2. This was the number of columns, minus the number of rows having leading entries (or m − rank
A). It is convenient to write solutions of AX = O as column vectors. In the last example, we could write ⎛ 35 ⎞  − 13  16 16 ⎜ 28 ⎟  ⎜− 16  + 20 16 ⎟ ⎜ ⎟ X = ⎜− 7  + 9 ⎟  ⎜ 16 16 ⎟ ⎝ ⎠   This formulation also makes it easy to display other information about solutions. In this example, we can also write ⎞ ⎛ ⎞ ⎛ −13 35 ⎜ 20⎟ ⎜−28⎟ ⎟ ⎜ ⎟ ⎜ ⎟ +  ⎜ 9⎟  −7 X=!⎜ ⎟ ⎜ ⎟ ⎜ ⎝ 0⎠ ⎝ 16⎠ 16 0 in which ! = /16 can be any number (since  can be any number), and  = /16 is also any number. This displays the solution as a linear combination of two linearly independent vectors. We will say more about the significance of this in the next section.

EXAMPLE 7.25

Solve the system −x2 + 2x3 + 4x4 = 0 −x3 + 3x4 = 0 2x1 + x2 + 3x3 + 7x4 = 0 6x1 + 2x2 + 10x3 + 28x4 = 0

7.5 Solution of Homogeneous Systems of Linear Equations Let

⎛

0 ⎜0 A=⎜ ⎝2 6 We find that

⎛

1 ⎜0 AR = ⎜ ⎝0 0

−1 0 1 2

0 1 0 0

⎞ 4 3⎟ ⎟ 7⎠ 28

2 −1 3 10

0 0 1 0

275

⎞ 13 −10⎟ ⎟ −3⎠ 0

From the first three rows of AR , read that x1 + 13x4 = 0 x2 − 10x4 = 0 x3 − 3x4 = 0 Thus the solution is given by x1 = −13 x2 = 10 x3 = 3 x4 =  in which  can be any number. We can write the solution as ⎛ ⎞ −13 ⎜ 10⎟ ⎟ X = ⎜ ⎝ 3⎠ 1 with  any number. In this example every solution is a constant multiple of one 4− vector. Note also that m − rank
A = 4 − 3 = 1. We will now firm up some of the ideas we have discussed informally, and then look at additional examples. First, everything we have done in this section has been based on the assertion that AX = 0 and AR X = O have the same solutions. We will prove this. THEOREM 7.16

Let A be an n × m matrix. Then the linear homogeneous systems AX = O and AR X = O have the same solutions. We know that there is an n × n matrix  such that A = AR . Further,  can be written as a product of elementary matrices E1 · · · Er . Suppose first that X = C is a solution of AX = O. Then AC = O, so

Proof


AC =
AC = AR C = O = O Then C is also a solution of AR X = O. Conversely, suppose K is a solution of AR X = O. Then AR K = O. We want to show that AK = O also. Because AR K = O, we have
AK = O , or
Er · · · E1 AK = O

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By Theorem 7.8, for each Ej , there is elementary matrix Ej∗ that reverses the effect of Ej . Then, from the last equation we have ∗ E1∗ E2∗ · · · Er−1 Er∗
Er Er−1 · · · E2 E1 AK = O ∗ But Er∗ Er = In , because Er∗ reverses the effect of Er . Similarly, Er−1 Er−1 = In , until finally the last equation becomes

E1∗
E1 AK = AK = O Thus K is a solution of AX = O and the proof is complete. The method for solving AX = O, which we illustrated above, is called the of the GaussJordan method, or complete pivoting. Here is an outline of the method. Keep in mind that, in a system AX = O, row k gives the coefficients of equation k, and column j contains the coefficients of xj as we look down the set of equations. Gauss-Jordan Method for Solving AX = O 1. Find AR . 2. Look down the columns of AR . If column j contains the leading entry of some row (so all other elements of this column are zero), then xj is said to be dependent. Determine all the dependent unknowns. The remaining unknowns (if any) are said to be independent. 3. Each nonzero row of AR represents an equation in the reduced system, having one dependent unknown (in the column having the leading entry 1) and all other unknowns in this equation (if any) independent. This enables us to write this dependent unknown in terms of the independent ones. 4. After step (3) is carried out for each nonzero row, we have each dependent unknown in terms of the independent ones. The independent unknowns can then be assigned any values, and these determine the dependent unknowns, solving the system. We can write the resulting solution as a linear combination of column solutions, one for each independent unknown. The resulting expression, containing an arbitrary constant for each independent unknown, is called the general solution of the system.

EXAMPLE 7.26

Solve the system −x1 + x3 + x4 + 2x5 = 0 x2 + 3x3 + 4x5 = 0 x1 + 2x2 + x3 + x4 + x5 = 0 −3x1 + x2 + 4x5 = 0 The matrix of coefficients is

⎛

−1 ⎜ 0 A=⎜ ⎝ 1 −3

0 1 2 1

1 3 1 0

1 0 1 0

⎞ 2 4⎟ ⎟ 1⎠ 4

7.5 Solution of Homogeneous Systems of Linear Equations We find that

⎛

1

⎜ ⎜0 AR = ⎜ ⎜0 ⎝ 0

0

0

0

1

0

0

0

1

0

0

0

1

− 98

277

⎞

5⎟ ⎟ 8⎟  9⎟ 8⎠ − 41

Because columns 1 through 4 of AR contain of leading entries of rows, x1 , x2 , x3 and x4 are dependent, while the remaining unknown, x5 is independent. The equations of the reduced system (which has the same solutions as the original system) are: 9 x 1 − x5 = 0 8 5 x 2 + x5 = 0 8 9 x3 + x 5 = 0 8 1 x4 − x5 = 0 4 We wrote these out for illustration, but in fact the solution can be read immediately from AR . We can choose x5 = , any number, and then 9 5 9 1 x1 =  x2 = −  x3 = −  x4 =  8 8 8 4 The dependent unknowns are given by AR in terms of the independent unknowns (only one in this case). We can write this solution more neatly as ⎛ ⎞ 9 ⎜−5⎟ ⎜ ⎟ ⎟ X=!⎜ ⎜−9⎟ ⎝ 2⎠ 8 in which ! = /8 can be any number. This is the general solution of AX = O. In this example, m − rank
A = 5 − 4 = 1

EXAMPLE 7.27

Consider the system 3x1 − 11x2 + 5x3 = 0 4x1 + x2 − 10x3 = 0 4x1 + 9x2 − 6x3 = 0 The matrix of coefficients is

⎛

3 A = ⎝4 4

−11 1 9

⎞ 5 −10⎠  −6

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Matrices and Systems of Linear Equations

The reduced matrix is ⎛

1 AR = ⎝0 0

0 1 0

⎞ 0 0 ⎠ = I3  1

The reduced system is just x1 = 0 x2 = 0 x3 = 0 This system has only the trivial solution, with each xj = 0. Notice that in this example there are no independent unknowns. If there were, we could assign them any values and have infinitely many solutions.

EXAMPLE 7.28

Consider the system 2x1 − 4x2 + x3 + x4 + 6x5 + 12x6 − 5x7 = 0 −4x1 + x2 + 6x3 + 3x4 + 10x5 − 9x6 + 8x7 = 0 7x1 + 2x2 + 4x3 − 8x4 + 6x5 − 5x6 + 15x7 = 0 2x1 + x2 + 6x3 + 3x4 − 4x5 − 2x6 − 21x7 = 0 The coefficient matrix is ⎛

2 ⎜−4 A=⎜ ⎝ 7 2

⎞ −4 1 1 6 12 −5 1 6 3 10 −9 8⎟ ⎟ 2 4 −8 6 −5 15⎠ 1 6 3 −4 −2 −21

We find that ⎛

1

0

0

0

⎜ ⎜0 AR = ⎜ ⎜0 ⎝

1

0

0

0

1

0

0

0

0

1

−7 3 −233 82 1379 738 −1895 738

7 6 −395 164 −995 1476 1043 1476

⎞

−29 6 −375 ⎟ ⎟ 164 ⎟  2161 ⎟ 1476 ⎠ −8773 1476

From this matrix we see that x1 , x2 , x3 , x4 are dependent, and x5 , x6 , x7 are independent. We read immediately from AR that 7 7 29 x1 = x5 − x6 + x7 3 6 6 233 395 375 x5 + x6 + x x2 = 82 164 164 7 1379 995 2161 x5 + x6 − x x3 = − 738 1476 1476 7 1895 1043 8773 x4 = x5 − x6 + x 738 1476 1476 7

7.5 Solution of Homogeneous Systems of Linear Equations

279

while x5 , x6 and x7 can (independent of each other) be assigned any numerical values. To make the solution look neater, write x5 = 1476, x6 = 1476 and x7 = 1476!, where ,  and ! are any numbers. Now the solution can be written x1 = 3444 − 1722 + 7134! x2 = 4194 + 3555 + 3375! x3 = −2758 + 995 − 2161! x4 = 3790 − 1043 + 8773! x5 = 1476 x6 = 1476 x7 = 1476! with ,  and ! any numbers. In column notation, ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ 7134 −1722 3444 ⎜ 3375⎟ ⎜ 3555⎟ ⎜ 4194⎟ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎜−2161⎟ ⎜ 995⎟ ⎜−2758⎟ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ X = ⎜ ⎜ 3790⎟ +  ⎜−1043⎟ + ! ⎜ 8773⎟  ⎜ ⎟ ⎜ ⎜ 1476⎟ 0⎟ 0⎟ ⎟ ⎜ ⎜ ⎟ ⎜ ⎝ ⎠ ⎝ ⎠ ⎝ 0⎠ 1476 0 1476 0 0 This is the general solution, being a linear combination of three linearly independent 7− vectors. In this example, m − rank
A = 7 − 4 = 3 In each of these examples, after we found the general solution, we noted that the number m−rank
A, the number of columns minus the rank of A, coincided with the number of linearly independent column solutions in the general solution (the number of arbitrary unknowns in the general solution). We will see shortly that this is always true. In the next section we will put the Gauss-Jordan method into a vector space context. This will result in an understanding of the algebraic structure of the solutions of a system AX = O, as well as practical criteria for determining when such a system has a nonzero solution.

SECTION 7.5

PROBLEMS

In each of Problems 1 through 12, find the general solution of the system and write it as a column matrix or sum of column matrices. 1.

x1 − 2x3 + x5 = 0 x3 + x4 − x5 = 0 6.

x1 − 2x5 = 0 7.

−3x3 + 2x4 + x5 = 0 3.

x 1 + x2 = 0 4.

−10x1 − x2 + 4x3 − x4 + x5 − x6 = 0 x2 − x3 + 3x4 = 0 2x1 − x2 + x5 = 0

−2x1 + x2 + 2x3 = 0 x 1 − x2 = 0

6x1 − x2 + x3 = 0 x1 − x4 + 2x5 = 0

−3x1 + x2 − x3 + x4 + x5 = 0 x2 + x3 + 4x5 = 0

x1 − x2 + 3x3 − x4 + 4x5 = 0 2x1 − 2x2 + x3 + x4 = 0

x1 + 2x2 − x3 + x4 = 0 x2 − x3 + x4 = 0

2.

5.

x2 − x4 + x6 = 0 8.

8x1 − 2x3 + x6 = 0 2x1 − x2 + 3x4 − x6 = 0

4x1 + x2 − 3x3 + x4 = 0

x2 + x3 − 2x5 − x6 = 0

2x1 − x3 = 0

x4 − 3x5 + 2x6 = 0

CHAPTER 7

280

Matrices and Systems of Linear Equations

x2 − 3x4 + x5 = 0

9.

11.

x1 − 2x2 + x5 − x6 + x7 = 0

2x1 − x2 + x4 = 0

x3 − x4 + x5 − 2x6 + 3x7 = 0

2x1 − 3x2 + 4x5 = 0

x1 − x5 + 2x6 = 0 2x1 − 3x4 + x5 = 0

10.

4x1 − 3x2 + x4 + x5 − 3x6 = 0

12.

2x2 − x6 + x7 − x8 = 0

2x2 + 4x4 − x5 − 6x6 = 0

x3 − 4x4 + x8 = 0

3x1 − 2x2 + 4x5 − x6 = 0

x2 − x3 + x4 = 0

2x1 + x2 − 3x3 + 4x4 = 0

7.6

2x1 − 4x5 + x7 + x8 = 0

x2 − x5 + x6 − x7 = 0

The Solution Space of AX = O Suppose A is an n×m matrix. We have been writing solutions of AX = O as column m−vectors. Now observe that the set of all solutions has the algebraic structure of a subspace of Rm . THEOREM 7.17

Let A be an n × m matrix. Then the set of solutions of the system AX = O is a subspace of Rm . Let S be the set of all solutions of AX = O. Then S is a set of vectors in Rm . Since AO = O, O is in S. Now suppose X1 and X2 are solutions, and  and  are real numbers. Then

Proof

A
X1 + X2  = AX1 + AX2 = O + O = O so X1 + X2 is also a solution, hence is in S. Therefore S is a subspace of Rm . We would like to know a basis for this solution space, because then every solution is a linear combination of the basis vectors. This is similar to finding a fundamental set of solutions for a linear homogeneous differential equation, because then every solution is a linear combination of these fundamental solutions. In examples in the preceding section, we were always able to write the general solution as a linear combination of m − rank
A linearly independent solution (vectors). This suggests that this number is the dimension of the solution space. To see why this is true in general, notice that we obtain a dependent xj corresponding to each row having a leading entry. Since only nonzero rows have leading entries, the number of dependent unknowns is the number of nonzero rows of AR . But then the number of independent unknowns is the total number of unknowns, m, minus the number of dependent unknowns (the number of nonzero rows of AR ). Since the number of nonzero rows of AR is the rank of A, then the general solution can always be written as a linear combination of m − rank
A independent solutions. Further, as a practical matter, solving the system AX = O by solving the system AR X = O automatically displays the general solution as a linear combination of this number of basis vectors for the solution space. We will summarize this discussion as a theorem.

7.6 The Solution Space of AX=O

281

THEOREM 7.18

Let A be n × m. Then the solution space of the system AX = O has dimension m − rank
A or, equivalently m − (number of nonzero rows in AR ).

EXAMPLE 7.29

Consider the system −4x1 + x2 + 3x3 − 10x4 + x5 = 0 2x1 + 8x2 − x3 − x4 + 3x5 = 0 −6x1 + x2 + x3 − 5x4 − 2x5 = 0 The matrix of coefficients is

⎛

−4 A=⎝ 2 −6 and we find that

⎛

1

⎜ A R = ⎝0 0

1 3 8 −1 1 1

0

0

1

0

0

1

⎞ −10 1 −1 3⎠  −5 −2 33 118 −32 59 −164 59

65 ⎞ 118 21 ⎟  59 ⎠ 56 59

Now A has m = 5 columns, and AR has 3 nonzero rows, so rank
A = 3 and the solution space of AX = O has dimension m − rank
A = 5 − 3 = 2 From the reduced system we read the solutions 65 21 33 32 −  x2 =  −  118 118 59 59 56 164 x3 =  −  x4 =  x5 =  59 59

x1 = −

in which  and  are any numbers. It is neater to replace  with 118! and  with 118 (which still can be any numbers) and write the general solution as ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ −65 −33 −33! − 65 ⎜ −42⎟ ⎜ 64⎟ ⎜ 64! − 42 ⎟ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ = ! ⎜ 328⎟ +  ⎜−112⎟  328! − 112 X=⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎝ ⎝ 118⎠ ⎠ ⎝ 0⎠ 118! 118 0 118 This displays the general solution (arbitrary element of the solution space) as a linear combination of two linearly independent vectors which form a basis for the dimension two solution space.

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We know that a system AX = O always has at least the zero (trivial) solution. This may be the only solution it has. Rank provides a useful criterion for determining when a system AX = O has a nontrivial solution. THEOREM 7.19

Let A be n × m. Then the system AX = O has a nontrivial solution if and only if m > rank
A This means that the system of homogeneous equations has a nontrivial solution exactly when the number of unknowns exceeds the rank of the coefficient matrix (the number of nonzero rows in the reduced matrix). We have seen that the dimension of the solution space is m − rank
A. There is a nontrivial solution if and only if this solution space has something in it besides the zero solution, and this occurs exactly when the dimension of this solution space is positive. But m − rank
A > 0 is equivalent to m > rank
A.

Proof

This theorem has important consequences. First, suppose the number of unknowns exceeds the number of equations. Then n < m. But rank
A ≤ n is always true, so in this case m − rank
A ≥ m − n > 0 and by Theorem 7.19, the system has nontrivial solutions. COROLLARY 7.2

A homogeneous system AX = O with more unknowns than equations always has a nontrivial solution. For another consequence of Theorem 7.19, suppose that A is square, so n = m. Now the dimension of the solution space of AX = O is n − rank
A. If this number is positive, the system has nontrivial solutions. If n − rank
A is not positive, then it must be zero, because rank
A ≤ n is always true. But n − rank
A = 0 corresponds to a solution space with only the zero solution. And it also corresponds, by Theorem 7.15, to A having the identity matrix as its reduced matrix. This means that a square system AX = O, having as the same number of unknowns as equations, has only the trivial solution exactly when the reduced form of A is the identity matrix. COROLLARY 7.3

Let A be an n × n matrix of real numbers. Then the system AX = O has only the trivial solution exactly when AR = In .

EXAMPLE 7.30

Consider the system −4x1 + x2 − 7x3 = 0 2x1 + 9x2 − 13x3 = 0 x1 + x2 + 10x3 = 0

7.7 Nonhomogeneous Systems of Linear Equations The matrix of coefficients is 3 × 3:

⎛

⎞ 1 −7 9 −13⎠  1 10

−4 A=⎝ 2 1 We find that

⎛

1 AR = ⎝0 0

283

0 1 0

⎞ 0 0 ⎠ = I3  1

This means that the system AX = O has only the solution x1 = x2 = x3 = 0. This makes sense in view of the fact that the system AX = O has the same solutions as the reduced system AR X = O, and when AR = I3 this reduced system is just X = O.

SECTION 7.6

PROBLEMS

1.–12. For n = 1 · · ·  12, use the solution of Problem n, Section 7.5, to determine the dimension of the solution space of the system of homogeneous equations.

14. Prove Corollary 7.2.

13. Can a system AX = O, in which there are at least

15. Prove Corollary 7.3.

7.7

as many equations as unknowns, have a nontrivial solution?

Nonhomogeneous Systems of Linear Equations We will now consider nonhomogeneous linear systems of equations: a11 x1 + a12 x2 + · · · + a1m xm = b1 a2 x1 + a22 x2 + · · · + a2m xm = b2   an1 x1 + an2 x2 + · · · + anm xm = bn  We can write this system in matrix form as AX = B, in which A = aij  is the n × m matrix of coefficients, ⎛ ⎞ ⎛ ⎞ b1 x1 ⎜b2 ⎟ ⎜ x2 ⎟ ⎜ ⎟ ⎜ ⎟ X = ⎜  ⎟ and B = ⎜  ⎟  ⎝  ⎠ ⎝  ⎠ xm bn This system has n equations in m unknowns. Of course, if each bj = 0 then this is a homogeneous system AX = O. A homogeneous system always has at least one a solution, the zero solution. A nonhomogeneous system need not have any solution at all.

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EXAMPLE 7.31

Consider the system 2x1 − 3x2 = 6 4x1 − 6x2 = 18 If there were a solution x1 =  x2 = , then from the first equation we would have 2 − 3 = 6. But then the second equation would give us 4−6 = 18 = 2
2−3 = 12, a contradiction. We therefore have an existence question to worry about with the nonhomogeneous system. Before treating this issue, we will ask: what must solutions of AX = B look like?

7.7.1

The Structure of Solutions of AX = B

We can take a cue from linear second order differential equations. There we saw that the every solution of y

+ py
+ qy = f
x is a sum of a solution the homogeneous equation y

+ py
+ qy = 0, and a particular solution of y

+ py
+ qy = f
x. We will show that the same idea holds true for linear algebraic systems of equations as well. THEOREM 7.20

Let Up be any solution of AX = B. Then every solution of AX = B is of the form Up + H, in which H is a solution of AX = O. Proof

Let W be any solution of AX = B. Since Up is also a solution of this system, then A
W − Up  = AW − AUp = B − B = O

Then W − Up is a solution of AX = O. Letting H = W − Up , then W = Up + H. Conversely, if W = Up + H, where H is a solution of AX = O, then AW = A
Up + H = AUp + AH = B + O = B so W is a solution of AX = B. This means that, if Up is any solution of AX = B, and H is the general solution of AX = O, then the expression Up + H contains all possible solutions of AX = B. For this reason we call such Up + H the general solution of AX = B, for any particular solution Up of AX = B.

EXAMPLE 7.32

Consider the system −x1 + x2 + 3x3 = −2 x2 + 2x3 = 4 Here

 −1 A= 0

1 3 1 2





and

 −2 B=  4

7.7 Nonhomogeneous Systems of Linear Equations

285

We find from methods of the preceding sections that the general solution of AX = O is ⎛ ⎞ 1  ⎝−2⎠  1 By a method we will describe shortly,

⎛ ⎞ 6 Up = ⎝ 4 ⎠ 0

is a particular solution of AX = B. Therefore every solution of AX = B is contained in the expression ⎛ ⎞ ⎛ ⎞ 1 6  ⎝−2⎠ + ⎝4⎠  1 0 in which  is any number. This is the general solution of the system AX = B.

7.7.2

Existence and Uniqueness of Solutions of AX = B

Now we know what to look for in solving AX = B. In this section we will develop criteria to determine when a solution Up exists, as well as a method that automatically produces the general solution in the form X = H + Up , where H is the general solution of AX = O

DEFINITION 7.14

Consistent System of Equations

A nonhomogeneous system AX = B is said to be consistent if there exists a solution. If there is no solution, the system is inconsistent.

The difference between a system AX = O and AX = B is B. For the homogeneous system, it is enough to specify the coefficient matrix A when working with the system. But for AX = B, we must incorporate B into our computations. For this reason, we introduce the  augmented matrix AB. If A is n × m, AB is the n ×
m + 1 matrix formed by adjoining B to A as a new last column. For example, if ⎛ ⎞ ⎛ ⎞ −3 2 6 1 5 A = ⎝ 0 3 3 −5⎠ and B = ⎝ 2⎠ 2 4 4 −6 −8 then

⎛

⎜−3  ⎜ AB = ⎜ 0 ⎝ 2

2

6

3

3

4

4

 1   −5   −6 

⎞ 5⎟ ⎟  2⎟ ⎠ −8

The column of dots does not count in the dimension of the matrix, and is simply a visual device to clarify that we are dealing with an augmented matrix giving both A and B for a system AX = B. If we just attached B as a last column without such an indicator, we might we dealing with a homogeneous system having 3 equations in 5 unknowns.

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Continuing with these matrices for the moment, reduce A to find AR : ⎞ ⎛ 1 1 0 0 3 ⎟ ⎜ AR = ⎝0 1 0 −3⎠  4 0 0 1 3  Next, reduce AB (ignore the dotted column in the row operations) to get ⎛ ⎜1  ⎜

ABR = ⎜0 ⎝ 0

0

0

1

0

0

1

   −3   4  3 1 3

⎞ − 163 ⎟ ⎟ 15 ⎟  4 ⎠ 37 − 12

Notice that  

ABR = AR C  If we reduce the augmented matrix AB, we obtain in the first m columns the reduced form of A, together with some new last column. The reason for this can be seen by reviewing how we reduce a matrix. Perform elementary row operations, beginning with the left-most column containing a leading entry, and work from left to right through the columns of the matrix. In  finding the reduced form of the augmented matrix AB, we deal with columns 1 · · · m, which  constitute A. The row operations used to reduce AB will, of course, operate on the elements of the last column as well, eventually resulting in what is called C in the last equation. We will state this result as a theorem.

THEOREM 7.21

Let A be n × m and let B be m × 1. Then for some m × 1 matrix C,  

ABR = AR C The reason this result is important is that the original system AX = B and the reduced system AR X = C have the same solutions (as in the homogeneous case, because the elementary row operations do not change the solutions of the system). But because of the special form of AR , the system AR X = C is either easy to solve by inspection, or to see that there is no solution.

EXAMPLE 7.33

Consider the system ⎛

−3 ⎝ 1 0

2 4 −2

⎞ ⎛ ⎞ 8 2 −6⎠ X = ⎝ 1⎠  −2 2

7.7 Nonhomogeneous Systems of Linear Equations We will reduce the augmented matrix ⎛ ⎜−3  ⎜

AB = ⎜ 1 ⎝ 0

2

2

4

−6

−2

2

⎞

     

One way to proceed is

287

8⎟ ⎟  1⎟ ⎠ −2 ⎛

⎞ 1⎟ ⎟ 2 2 8⎟ ⎠ 0 −2 2 −2 ⎞ ⎛   1 1 4 −6 ⎜ ⎟ ⎜ ⎟ → add 3(row 1) to row 2 → ⎜0 14 −16  11⎟ ⎝ ⎠  0 −2 2  −2 ⎛ ⎞  4 −6  1⎟ ⎜1 1 ⎜ ⎟  8 11  →
row 2 → ⎜0 ⎟ 1 −7  14 ⎠ ⎝ 14  0 −2 2  −2 ⎛ ⎞  10 15 1 0 −  − 7 7 ⎟ ⎜ ⎜ ⎟ 11 ⎟ → −4(row 2) to row 1 2(row 2) to row 3 → ⎜0 1 − 8  7 14 ⎠ ⎝  0 0 − 27  − 37 ⎛ ⎞  10 15 1 0 −  − 7 7 ⎟ ⎜ 7 ⎜ ⎟ 11 ⎟ → − (row 3) → ⎜0 1 − 8  7 14 ⎠ ⎝ 2  3 0 0 1  2 ⎛ ⎞  ⎜1 0 0  0 ⎟ ⎜ ⎟ 10 8  → (row 3) to row 1, (row 3) to row 2 → ⎜ 0 1 0  25 ⎟ ⎜ ⎟ 7 7 ⎝ ⎠  0 0 1  23 ⎜ 1  ⎜

AB → interchange rows 1 and 2 → ⎜−3 ⎝

4

−6

     

As can be seen in this process, we actually arrived at AR in the first three rows and columns, and whatever ends up in the last column is what we call C:  

ABR = AR C  Notice that the reduced augmented matrix is I3 C and represents the reduced system I3 X = C. This is the system ⎛ ⎞ ⎛ ⎞ 0 1 0 0 5⎟ ⎝0 1 0 ⎠ X = ⎜ ⎝2⎠ 3 0 0 1 2

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Matrices and Systems of Linear Equations

 which we solve by inspection to get x1 = 0 x2 = 5/2 x3 = 3/2. Thus reducing AB immediately yields the solution ⎛ ⎞ 0 ⎜5⎟ Up = ⎝ 2 ⎠ 3 2

of the original system AX = B. Because AR = I3 , Corollary 7.3 tells us that the homogeneous system AX = O has only the trivial solution, and therefore H = O in Theorem 7.20 and Up is the unique solution of AX = B.

EXAMPLE 7.34

The system 2x1 − 3x2 = 6 4x1 − 6x2 = 18 is inconsistent, as we saw in Example 7.31. We will put the fact that this system has no solution into the context of the current discussion. Write the augmented matrix ⎛ ⎞   2 −3  6 ⎠

AB = ⎝  4 −6  18 Reduce this matrix. We find that ⎛  1

ABR = ⎝ 0

− 23 0

   

⎞ 0⎠



1

From this we immediately read the reduced system AR X = C :     0 1 − 23 X=  AR X = 1 0 0 This system has the same solutions as the original system. But the second equation of the reduced system is 0x1 + 0x2 = 1 which has no solution. Therefore AX = B has no solution either. In this example, the reduced system has an impossible equation because AR has a zero second row, while the second row of ABR has a nonzero element in the augmented column. Whenever this happens, we obtain an equation having all zero coefficients of the unknowns, but equal to a nonzero number. In such a case the reduced system AR X = C, hence the original system AX = B, can have no solution. The key to recognizing when this will occur is that it happens when the rank of AR (its number of nonzero rows) is less than the rank of AB.

7.7 Nonhomogeneous Systems of Linear Equations

289

THEOREM 7.22

 The nonhomogeneous system AX = B has a solution if and only if A and AB have the same rank.  Proof Let A be n × m. Suppose  first that rank
A = rank
AB = r. By Theorems 7.12 and 7.14, the column space of AB has dimension r.  Certainly r cannot exceed the number of columns of A, soB, which is column m + 1 of AB, must be a linear combination of the first m columns of AB, which form A. This means that, for some numbers 1      m , ⎛

⎞ ⎛ ⎞ ⎛ ⎞ a11 a12 a1m ⎜a21 ⎟ ⎜a22 ⎟ ⎜a2m ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ B = 1 ⎜  ⎟ + 2 ⎜  ⎟ + · · · + m ⎜  ⎟ ⎝  ⎠ ⎝  ⎠ ⎝  ⎠ an1

⎛

an2

⎞

anm ⎛

⎞ 1 1 a11 + 2 a12 + · · · + m a1m ⎜ 1 a21 + 2 a22 + · · · + m a2m ⎟ ⎜ 2 ⎟ ⎜ ⎟ ⎜ ⎟ =⎜ ⎟ = A ⎜  ⎟   ⎝  ⎠ ⎝ ⎠  1 an1 + 2 an2 + · · · + m anm m ⎛

⎞ 1 ⎜ 2 ⎟ ⎜ ⎟ But then ⎜  ⎟ is a solution of AX = B. ⎝  ⎠ m

⎛

⎞ 1 ⎜ 2 ⎟ ⎜ ⎟ Conversely, suppose AX = B has a solution ⎜  ⎟. Then ⎝  ⎠ m ⎞ ⎛ ⎞ 1 a11 + 2 a12 + · · · + m a1m 1 ⎜ 2 ⎟ ⎜ 1 a21 + 2 a22 + · · · + m a2m ⎟ ⎟ ⎜ ⎟ ⎜ B = A⎜  ⎟ = ⎜ ⎟   ⎠ ⎝  ⎠ ⎝  ⎛

m 1 an1 + 2 an2 + · · · + m anm ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ a11 a12 a1m ⎜a21 ⎟ ⎜a22 ⎟ ⎜a2m ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ = 1 ⎜  ⎟ +  2 ⎜  ⎟ + · · · +  m ⎜  ⎟  ⎝  ⎠ ⎝  ⎠ ⎝  ⎠ an1

an2

anm

n Then B is a linear combination of the columns of A, thought  of as vectors in R . But then the column space of A is the same as the column space of AB. Then

rank
A = dimension of the column space of A   = dimension of the column space of AB = rank AB and the proof is complete.

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EXAMPLE 7.35

Solve the system x1 − x2 + 2x3 = 3 −4x1 + x2 + 7x3 = −5 −2x1 − x2 + 11x3 = 14 The augmented matrix is

⎛

 ⎜ 1 −1 2   

AB = ⎜ ⎝−4 1 7  −2 −1 11 · · · When we reduce this matrix we obtain

⎞ 3⎟ ⎟ −5⎠ 14

⎛

⎜1   ⎜

ABR = AR C = ⎜0 ⎝ 0

0

−3

1

−5

0

0

⎞ 0⎟ ⎟  0⎟ ⎠ 1

     

The first three columns of this reduced matrix make up AR . But  rank
A = 2 and rank
ABR  = 3 so this system has no solution. The last equation of the reduced system is 0x1 + 0x2 + 0x3 = 1 which can have no solution.

EXAMPLE 7.36

Solve x1 − x3 + 2x4 + x5 + 6x6 = −3 x2 + x3 + 3x4 + 2x5 + 4x6 = 1 x1 − 4x2 + 3x3 + x4 + 2x6 = 0 The augmented matrix is

⎛

⎜1  ⎜

AB = ⎜0 ⎝ 1 Reduce this to get

0

−1

2

1

6

1

1

3

2

4

−4

3

1

0

2

⎛ ⎜1  ⎜

ABR = ⎜0 ⎝ 0

0

0

27 8

15 8

60 8

1

0

13 8

9 8

20 8

0

1

11 8

7 8

12 8

     

     

⎞ −3⎟ ⎟  1⎟ ⎠ 0 ⎞ − 178 ⎟ ⎟ 1⎟ 8⎠ 7 8

7.7 Nonhomogeneous Systems of Linear Equations

291

 The first six columns of this matrix form AR and we read that rank
A = 3 = rank
ABR . From ABR , identify x1  x2  x3 as dependent and x4  x5  x6 as independent. The number of independent unknowns is m − rank
A = 6 − 3 = 3, and this is the dimension of the solution space of AX = O. From the reduced augmented matrix, the first equation of the reduced system is x1 +

27 15 60 17 x4 + x5 + x6 = −  8 8 8 8

so 27 15 60 17 x4 − x5 − x6 −  8 8 8 8 We will not write out all of the equations of the reduced system. The point is that we can read  directly from AB that x1 = −

R

x2 = −

13 9 20 1 x4 − x5 − x6 + 8 8 8 8

and 11 7 12 7 x4 − x5 − x6 +  8 8 8 8 while x4  x5  x6 can be assigned any numerical values. We can write this solution as ⎛ 27 ⎞ − 8 x4 − 158 x5 − 608 x6 − 178 ⎜ 13 ⎟ ⎜ − 8 x4 − 98 x5 − 208 x6 + 18 ⎟ ⎜ ⎟ ⎜ − 11 x − 7 x − 12 x + 7 ⎟ ⎜ 4 5 6 8 8 8 ⎟ X=⎜ 8 ⎟ ⎜ ⎟ x4 ⎜ ⎟ ⎝ ⎠ x5 x6 x3 = −

If we let x4 = 8, x5 = 8, and x6 = 8!, with ,  and ! any numbers, then the general solution is ⎛ ⎞ ⎛ 17 ⎞ ⎛ ⎞ ⎛ ⎞ −8 −60 −15 −27 ⎜−20⎟ ⎜ 1 ⎟ ⎜ −9⎟ ⎜−13⎟ ⎜ ⎟ ⎜ 8⎟ ⎜ ⎟ ⎜ ⎟ ⎜−12⎟ ⎜ 7 ⎟ ⎜ −7⎟ ⎜−11⎟ ⎜ ⎟ ⎜ 8⎟ ⎜ ⎟ ⎜ ⎟ X = ⎜ ⎟ +  ⎜ 0⎟ + ! ⎜ 0⎟ + ⎜ 0⎟  ⎜ ⎟ ⎜ ⎜ ⎟ ⎜ 8⎟ ⎟ ⎝ 0⎠ ⎝ 0⎠ ⎝ 8⎠ ⎝ 0⎠ 8 0 0 0 This is in the form H + Up , with H the general solution of AX = O and Up a particular solution of AX = B. Since the general solution is of the form X = H + Up , with H the general solution of AX = O, the only way AX = B can have a unique solution is if H = O, that is, the homogeneous system must have only the trivial solution. But, for a system with the same number of unknowns as equations, this can occur only if AR is the identity matrix. THEOREM 7.23

Let A be n × n. Then the nonhomogeneous system AX = B has a unique solution if and only if AR = In . This, in turn, occurs exactly when rank
A = n.

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EXAMPLE 7.37

Consider the system

⎛

2 ⎝−5 1 The augmented matrix is

1 1 1

⎞ ⎛ ⎞ −11 −6 9⎠ X = ⎝ 12⎠  14 −5

⎞  ⎜ 2 1 −11  −6⎟  ⎟ ⎜ 

AB = ⎜−5 1 9  12⎟ ⎠ ⎝  1 1 14  −5 ⎛

and we find that

⎛ ⎜1  ⎜

ABR = ⎜0 ⎝

0

0

1

0

0

0

1

     

⎞ − 86 31 ⎟ ⎟ 191 ⎟  − 155 ⎠ 11 − 155

The first three columns tell us that AR = I3 . The homogeneous system AX = O has only the  trivial solution. Then AX = B has a unique solution, which we read from ABR : ⎛ −86 ⎞ 31

⎟ ⎜  X = ⎝ −191 155 ⎠ −11 155

Note that rank
A = 3 and the dimension of the solution space AX = O is n − rank
A = 3 − 3 = 0, consistent with this solution space having no elements except the zero vector.

SECTION 7.7

PROBLEMS

In each of Problems 1 through 14, find the general solution of the system or show that the system has no solution. 1.

4.

3x1 − 2x2 + x3 = 6

x1 − 4x2 = 3

x1 + 10x2 − x3 = 2 −3x1 − 2x2 + x3 = 0 2.

5.

x2 + x3 − 6x4 + x6 = −9 x1 − x2 + x6 = 0

2x1 − 3x2 + x4 = 16 2x1 − 3x2 + x4 − x6 = 0

3x2 − 4x4 = 10 x1 − 3x2 + 4x5 − x6 = 8

4x1 − 2x2 + 3x3 + 10x4 = 1 x1 − 3x4 = 8

3.

2x1 − 3x2 = 1 −x1 + 3x2 = 0

6.

2x1 − 3x2 + x4 = 1

3x1 − 2x3 + x5 = 1

3x2 + x3 − x4 = 0

x2 − x4 + 6x6 = 3

2x1 − 3x2 + 10x3 = 0

7.8 Matrix Inverses 7.

8x2 − 4x3 + 10x6 = 1

12.

−4x1 + 5x2 − 6x3 = 2

x3 + x5 − x6 = 2

2x1 − 6x2 + x3 = −5

x4 − 3x5 + 2x6 = 0 2x1 − 3x3 = 1

8.

−6x1 + 16x2 − 11x3 = 1 13.

x1 − x2 + x3 = 1 9. 10. 11.

4x1 + 3x2 = 4 7x1 − 3x2 + 4x3 = −7

−2x1 + x2 + 7x3 = 4 14.

−6x1 + 2x2 − x3 + x4 = 0 x1 + 4x2 − x4 = −5 x1 + x2 + x3 − 7x4 = 0

15. Let A be an n × m matrix with rank r. Prove that the reduced system AR X = B has a solution if and only if br+1 = · · · = bn = 0.

2x1 + x2 − x3 + 4x4 = 6 x2 − 3x4 = −5

7.8

4x1 − x2 + 4x3 = 1 x1 + x2 − 5x3 = 0

2x1 − 4x2 + x3 = 2 14x3 − 3x5 + x7 = 2 x1 + x2 + x3 − x4 + x6 = −4 3x1 − 2x2 = −1

293

Matrix Inverses

DEFINITION 7.15

Matrix Inverse

Let A be an n × n matrix. Then B is an inverse of A if AB = BA = In 

In this definition B must also be n × n because both AB and BA must be defined. Further, if B is an inverse of A, then A is also an inverse of B. It is easy to find nonzero square matrices that have no inverse. For example, let   1 0 A=  2 0   a b If B is an inverse of A, say B = , then we must have c d        1 0 a b a b 1 0 AB = = =  2 0 c d 2a 2b 0 1 But then a = 1 b = 0 2a = 0

and

b=1

which are impossible conditions. On the other hand, some matrices do have inverses. For example,   4    4    − 17 − 17 2 1 2 1 1 0 7 7 = =  2 2 1 4 1 4 0 1 − 17 − 17 7 7

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DEFINITION 7.16

Nonsingular and Singular Matrices

A square matrix is said to be nonsingular if it has an inverse. If it has no inverse, the matrix is called singular.

If a matrix has an inverse, then it can have only one. THEOREM 7.24

Uniqueness of Inverses

Let B and C be inverses of A. Then B = C. Proof

Write B = BIn = B
AC =
BAC = In C = C

In view of this we will denote the inverse of A as A−1 . Here are properties of inverse matrices. In proving parts of the theorem, we repeatedly employ the strategy that, if AB = BA = In , then B must be the inverse of A. THEOREM 7.25

1. In is nonsingular and In−1 = In . 2. If A and B are nonsingular n × n matrices, then AB is nonsingular and
AB−1 = B−1 A−1  3. If A is nonsingular, so is A−1 , and
A−1 −1 = A 4. If A is nonsingular, so is At , and
At −1 =
A−1 t  5. If A and B are n × n and either is singular, then AB and BA are both singular. Proof

For (2), compute
AB
B−1 A−1  = A
BB−1 A−1 = AA−1 = In 

Similarly,
B−1 A−1 
AB = In . Therefore
AB−1 = B−1 A−1 . For (4), use Theorem 7.3(3) to write
At 
A−1 t =
A−1 At =
In t = In  Similarly,
A−1 t
At  =
AA−1 t = In  Therefore
At −1 =
A−1 t . We will be able to give a very short proof of (5) when we have developed determinants. We saw before that not every matrix has an inverse. How can we tell whether a matrix is singular or nonsingular? The following theorem gives a reasonable test.

7.8 Matrix Inverses

295

THEOREM 7.26

An n × n matrix A is nonsingular if and only if AR = In . Alternatively, an n × n matrix is nonsingular if and only if its rank is n. The proof consists of understanding a relationship between a matrix having an inverse, and its reduced form being the identity matrix. The key lies in noticing that we can form the columns of a matrix product AB by multiplying, in turn, A by each column of B: ⎛ ⎞ b1j ⎜b2j ⎟ ⎜ ⎟ column j of AB = A(column j of B = A ⎜  ⎟  ⎝  ⎠ bnj We will build an inverse for A a column at a time. To have AB = In , we must be able to choose the columns of B so that ⎛ ⎞ 0 ⎛ ⎞ ⎜0⎟ b1j ⎜ ⎟ ⎜  ⎟ ⎜b2j ⎟ ⎜⎟ ⎜ ⎟ ⎟ column j of AB = A ⎜  ⎟ = column j of In = ⎜ (7.4) ⎜1⎟  ⎝  ⎠ ⎜ ⎟ ⎜⎟ bnj ⎝  ⎠ 0

Proof

with a 1 in the j th place and zeros elsewhere. Suppose now that AR = In . Then, by Theorem 7.23, the system (7.4) has a unique solution for each j = 1 · · ·  n. These solutions form the columns of a matrix B such that AB = In , and then B = A−1 . (Actually we must show that AB = In also, but we leave this as an exercise). Conversely, suppose A is nonsingular. Then system (7.4) has a unique solution for j = 1     n, because these solutions are the columns of A−1 . Then, by Theorem 7.23, AR = In .

7.8.1

A Method for Finding A−1

We know some computational rules for working with matrix inverses, as well as a criterion for a matrix to have an inverse. Now we want an efficient way of computing A−1 from A. Theorem 7.26 suggests a strategy. We know that, in any event, there is an n × n matrix  such that A = AR .  is a product of elementary matrices representing the elementary row operations used to reduce A. Previously we found  by adjoining In to the left of A to form   an n × 2n matrix In A. Reduce A, performing the elementary row operations on all of In A  to eventually arrive at In . This produces  such that A = AR  If AR = In then  = A−1 . If AR = In , then A has no inverse.

EXAMPLE 7.38

Let

 A=

5 6

 −1  8

We want to know if A is nonsingular and, if it is, produce its inverse.

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Matrices and Systems of Linear Equations

Form

⎛  1

I2 A = ⎝ 0

0 1

   

⎞ 5 −1⎠  6 8

Reduce A (the last two columns), carrying out the same ⎛ 1  1 0

I2 A → (row 1) → ⎝ 5 5 0 1 ⎛ 1 5

operations on the first two columns: ⎞   1 − 15 ⎠   6 8    

0

→ −6(row 1) + (row 2) → ⎝ − 65 1 ⎛ 1 5 0 → (row 2) → ⎝ 5 46 − 466 465 ⎛ 8 1 1 → (row 2) + (row 1) → ⎝ 46 46 5 − 466 465

⎞ − 15 ⎠

1

46 5

0        

1

⎞ − 15 ⎠

0

1 ⎞

1

0⎠

0

1



In the last two columns we read AR = I2 . This means that A is nonsingular. From the first two columns,   1 8 1  A−1 = 46 −6 5

EXAMPLE 7.39

Let



−3 A= 4 Perform a reduction:

 21  −28

⎛   1 0 

I2 A = ⎝  0 1  ⎛ 1 −1 0 → − (row 1) → ⎝ 3 3 0 1 ⎛ −1 0 −4(row 1) + (row 2) → ⎝ 3 4 1 3

−3 4        

⎞ 21⎠ −28

⎞ −7⎠ 4 −28 ⎞ 1 −7⎠  0 0 1

We read AR from the last two columns, which form a 2 × 2 reduced matrix. Since this is not I2 , A is singular and has no inverse. Here is how inverses relate to the solution of systems of linear equations in which the number of unknowns equals the number of equations.

7.8 Matrix Inverses

297

THEOREM 7.27

Let A be an n × n matrix. 1. A homogeneous system AX = O has a nontrivial solution if and only if A is singular. 2. A nonhomogeneous system AX = B has a solution if and only if A is nonsingular. In this case the unique solution is X = A−1 B. For a homogeneous system AX = O, if A were nonsingular then we could multiply the equation on the left by A−1 to get X = A−1 O = O Thus in the nonsingular case, a homogeneous system can have only a trivial solution. In the singular case, we know that rank
A < n, so the solution space has positive dimension n − rank
A and therefore has nontrivial solutions in it. In the nonsingular case, we can multiply a nonhomogeneous equation AX = B on the left by A−1 to get the unique solution X = A−1 B However, if A is singular, then rank
A < n, and then Theorem 7.22 tells us that the system AX = B can have no solution.

EXAMPLE 7.40

Consider the nonhomogeneous system 2x1 − x2 + 3x3 = 4 x1 + 9x2 − 2x3 = −8 4x1 − 8x2 + 11x3 = 15 The matrix of coefficients is

⎛

2 A = ⎝1 4 and we find that

⎞ −1 3 9 −2⎠ −8 11

⎛

⎞ 83 −13 −25 1 ⎝ −19 10 7⎠  A−1 = 53 −44 12 19

The unique solution of this system is

⎛ ⎞⎛ ⎞ 83 −13 −25 4 1 ⎝−19 10 7⎠ ⎝−8⎠ X = A−1 B = 53 −44 12 19 15 ⎛ 61 ⎞ 53

⎟ ⎜ = ⎝− 51  53 ⎠ 13 53

298

CHAPTER 7

SECTION 7.8

Matrices and Systems of Linear Equations

PROBLEMS

In each of Problems 1 through 10, find the inverse of the matrix or show that the matrix is singular.   −1 2 1. 2 1   12 3 2. 4 1   −5 2 3. 1 2   −1 0 4. 4 4   6 2 5. 3 3 ⎛ ⎞ 1 1 −3 ⎜ ⎟ 1⎠ 6. ⎝2 16 0 0 4 ⎞ ⎛ −3 4 1 ⎟ ⎜ 7. ⎝ 1 2 0⎠ 1 1 3 ⎞ ⎛ −2 1 −5 ⎟ ⎜ 1 4⎠ 8. ⎝ 1 0 3 3 ⎞ ⎛ −2 1 1 ⎟ ⎜ 9. ⎝ 0 1 1⎠ −3 0 6

⎛

12 ⎜ 10. ⎝−3 0

1 2 9

⎞ 14 ⎟ 0⎠ 14

In each of Problems 11 through 15, find the unique solution of the system, using Theorem 7.27(2). 11.

x1 − x2 + 3x3 − x4 = 1 x2 − 3x3 + 5x4 = 2 x1 − x3 + x4 = 0 x1 + 2x3 − x4 = −5

12.

8x1 − x2 − x3 = 4 x1 + 2x2 − 3x3 = 0 2x1 − x2 + 4x3 = 5

13.

2x1 − 6x2 + 3x3 = −4 −x1 + x2 + x3 = 5 2x1 + 6x2 − 5x3 = 8

14.

12x1 + x2 − 3x3 = 4 x1 − x2 + 3x3 = −5 −2x1 + x2 + x3 = 0

15.

4x1 + 6x2 − 3x3 = 0 2x1 + 3x2 − 4x3 = 0 x1 − x2 + 3x3 = −7

16. Let A be nonsingular. Prove that, for any positive integer k, Ak is nonsingular, and
Ak −1 =
A−1 k . 17. Let A, B and C be n × n real matrices. Suppose BA = AC = In . Prove that B = C.

CHAPTER

8

PERMUTAIONS DEFINITION OF THE DETERMINANT PROPERTIES EVALUATION OF DETERMINANTS ELEMENTARY ROW AND COLUMN OPERATIONS COFACTOR EXPANSIONS DETERMINANTS OF TRIANGULAR

Determinants

If A is a square matrix, the determinant of A is a sum of products of elements of A, formed according to a procedure we will now describe. First we need some information about permutations.

8.1

Permutations If n is a positive integer, a permutation of order n is an arrangement of the integers 1     n in any order. For example, suppose p is a permutation that reorders the integers 1     6 as 3 1 4 5 2 6 Then p
1 = 3

p
2 = 1

p
3 = 4

p
4 = 5

p
5 = 2

p
6 = 6

with p
j the number the permutation has put in place j. For small n it is possible to list all permutations on 1     n. Here is a short list: For n = 2 there are two permutations on the integers 1 2, one leaving them in place and the second interchanging them: 1 2 2 1 299

300

CHAPTER 8

Determinants

For n = 3 there are six permutations on 1 2 3, and they are

1 2 3 1 3 2 2 1 3 2 3 1 3 1 2 3 2 1

For n = 4 there are twenty four permutations on 1 2 3 4: 1 2 3 4 1 2 4 3 1 3 2 4 1 3 4 2 1 4 2 3 1 4 3 2 2 1 3 4 2 1 4 3 2 3 1 4 2 3 4 1 2 4 1 3 2 4 3 1 3 1 2 4 3 1 4 2 3 2 1 4 3 2 4 1 3 4 1 2 3 4 2 1 4 1 2 3 4 1 3 2 4 2 1 3 4 2 3 1 4 3 1 2 4 3 2 1 An examination of this list of permutations suggests a systematic approach by which they were all listed, and such an approach will work in theory for higher n. However, we can also observe that the number of permutations on 1     n increases rapidly with n. There are n! = 1 · 2 · · · · · n permutations on 1     n. This fact is not difficult to derive. Imagine a row of n boxes, and start putting the integers from 1 to n into the boxes, one to each box. There are n choices for a number to put into the first box, n − 1 choices for the second, n − 2 for the third, and so on until there is only one left to put in the last box. There is a total of n
n − 1
n − 2 · · · 1 = n! ways to do this, hence n! permutations on n objects. A permutation is characterized as even or odd, according to a rule we will now illustrate. Consider the permutation 2 5 1 4 3 on the integers 1     5. For each number k in the list, count the number of integers to its right that are smaller than k. In this way form a list k 2 5 1 4 3

number of integers smaller than k to the right of k 1 3 0 1 0

Sum the integers in the right column to get 5, which is odd. We therefore call this permutation odd. As an example of an even permutation, consider 2 1 5 4 3

8.2 Definition of the Determinant

301

Now the list is k 2 1 5 4 3

number of integers smaller than k to the right of k 1 0 2 1 0

and the integers in the right column sum to 4, an even number. This permutation is even. If p is a permutation, let  0 if p is even sgn
p =  1 if p is odd

SECTION 8.1

PROBLEMS

1. The six permutations of 1 2 3 are given in the discussion. Which of these permutations are even and which are odd?

3. Show that half of the permutations on 1 2     n are even, and the other half are odd.

2. The 24 permutations of 1 2 3 4 are given in the discussion. Which of these are even and which are odd?

8.2

Definition of the Determinant Let A = aij  be an n × n matrix, with numbers or functions as elements.

DEFINITION 8.1

The determinant of A, denoted det
A, is the sum of all products
−1sgn
p a1p
1 a2p
2 · · · anp
n  taken over all permutations p on 1     n. This sum is denoted 
−1sgn
p a1p
1 a2p
2 · · · anp
n 

(8.1)

p

Each term in the defining sum (8.1) contains exactly one element from each row and from each column, chosen according to the indices j p
j determined by the permutation. Each product in the sum is multiplied by 1 if the permutation is even, and by −1 if p is odd.

302

CHAPTER 8

Determinants

Since there are n! permutations on 1     n, this sum involves n! terms and is therefore quite daunting for, say n ≥ 4. We will examine the small cases n = 2 and n = 3 and then look for ways of evaluating det
A for larger n. In the case n = 2,   a11 a12  A= a21 a22 We have seen that there are 2 permutations on 1 2, namely p  1 2 which is an even permutation, and q  2 1 which is odd. Then det
A =
−1sgn
p a1p
1 a2p
2 +
−1sgn
q a1q
1 a2q
2 =
−10 a11 a22 +
−11 a12 a21 = a11 a22 − a12 a21  This rule for evaluating det
A holds for any In the case n = 3, ⎛ a11 A = ⎝ a21 a31

2 × 2 matrix. a12 a22 a32

⎞ a13 a23 ⎠  a33

The permutations of 1 2 3 are p1  1 2 3 p2  1 3 2 p3  2 1 3 p4  2 3 1 p5  3 1 2 p6  3 2 1 It is routine to check that p1 , p5 and p6 are even, and p2 , p3 and p4 are odd. Then det
A =
−1sgn
p1  a1p1
1 a2p1
2 a3p1
3 +
−1sgn
p2  a1p2
1 a2p2
2 a3p2
3 +
−1sgn
p3  a1p3
1 a2p3
2 a3p3
3 +
−1sgn
p4  a1p4
1 a2p4
2 a3p4
3 +
−1sgn
p5  a1p5
1 a2p5
2 a3p5
3 +
−1sgn
p6  a1p6
1 a2p6
2 a3p6
3 = a11 a22 a33 − a11 a23 a32 − a12 a21 a33 + a12 a23 a31 + a13 a21 a32 − a13 a22 a31  If A is 4 × 4, then evaluation of det
A by direct recourse to the definition will involve 24 terms, as well as explicitly listing all 24 permutations on 1 2 3 4. This is not practical. We will therefore develop some properties of determinants which will make their evaluation more efficient.

8.3 Properties of Determinants

SECTION 8.2

PROBLEMS ⎛

6 3. A = ⎝ 2 0

In Problems 1 through 4, use the formula for det
A in the 3 × 3 case to evaluate the determinant of the given matrix. ⎛

1 1. A = ⎝ 1 0 ⎛

−1 2. A = ⎝ 2 1

8.3

⎞ 0 −1 ⎠ 1

6 2 1 3 2 1

303

⎛

−4 4. A = ⎝ 0 0

⎞ 1 0 ⎠ 4

−3 1 1 0 1 0

⎞ 5 4 ⎠ −4 ⎞ 1 1 ⎠ 0

5. The permutations on 1 2 3 4 were listed in Section 8.1. Use this list to write a formula for det
A when A is 4 × 4.

Properties of Determinants We will develop some of the properties of determinants that are used in evaluating them and deriving some of their properties. There are effective computer routines for evaluating quite large determinants, but these are also based on the properties we will display. First, it is standard to use vertical lines to denote determinants, so we will often write det
A = A  This should not be confused with absolute value. If A has numerical elements, then A is a number and can be positive, negative or zero. Throughout the rest of this chapter let A and B be n × n matrices. Our first result says that a matrix having a zero row has a zero determinant. THEOREM 8.1

If A has a zero row, then A = 0. This is easy to see from the defining sum (8.1). Suppose, for some i , each aij = 0. Each term of the sum (8.1) contains a factor aipj
i from row i, hence each term in the sum is zero. Next, we claim that multiplying a row of a matrix by a scalar  has the effect of multiplying the determinant of the matrix by . THEOREM 8.2

Let B be formed from A by multiplying row k by a scalar . Then B =  A  The effect of multiplying row k of A by  is to replace each akj by akj . Then bij = aij for i = k, and bkj = akj , so  A =
−1sgn
p a1p
1 a2p
2 · · · akp
k · · · anp
n p

304

CHAPTER 8

Determinants

and B =




−1sgn
p b1p
1 b2p
2 · · · bkp
k · · · bnp
n

p

=

 p

=


−1sgn
p a1p
1 a2p
2 · · ·
akp
k  · · · anp
n




−1sgn
p a1p
1 a2p
2 · · · akp
k · · · anp
n =  A 

p

The next result states that the interchange of two rows in a matrix causes a sign change in the determinant. THEOREM 8.3

Let B be formed from A by interchanging two rows. Then A = − B  A proof of this involves a close examination of the effect of a row interchange on the terms of the sum (8.1), and we will not go through these details. The result is easy to see in the case of 2 × 2 determinants. Let     a11 a12 a21 a22 and B =  A= a21 a22 a11 a12 Then A = a11 a22 − a12 a21 and B = a21 a12 − a22 a11 = − A  This result has two important consequences. The first is that the determinant of a matrix with two identical rows must be zero. COROLLARY 8.1

If two rows of A are the same, then A = 0. The reason for this is that, if we form B from A by interchanging the identical rows, then B = A, so B = A. But by Theorem 8.3, B = − A, so A = 0. COROLLARY 8.2

If for some scalar , row k of A is  times row i, then A = 0. To see this, consider two cases. First, if  = 0, then row k of A is a zero row, so A = 0. If  = 0, then we can multiply row k of A by 1/ to obtain a matrix B having rows i and k the same. Then B = 0. But B =
1/ A by Theorem 8.2, so A = 0. Next, we claim that the determinant of a product is the product of the determinants.

8.3 Properties of Determinants

305

THEOREM 8.4

Let A and B be n × n matrices. Then, AB = A B. Obviously this extends to a product involving any finite number of n × n matrices. The theorem enables us to evaluate the determinant of such a product without carrying out the matrix multiplications of all the factors. We will illustrate the theorem when we have efficient ways of evaluating determinants. The following theorem gives the determinant of a matrix that is written as a sum of matrices in a special way.

THEOREM 8.5

Suppose each element of row k of A is written as a sum kj + kj . Form two matrices from A. The first, A1 , is identical to A except the elements of row k are kj . The second, A2 , is identical to A except the elements of row k are kj . Then A = A1  + A2   If we display the elements of these matrices, the conclusion states that   a11        k1 + k1        an1

···  

···  

a1n  

· · · kj + kj    

···  

kn + kn  

···

···

ann

a1j  

anj

     a11 · · · a1j · · · a1n                        =  k1 · · · kj · · · kn                         an1 · · · anj · · · ann    a11 · · · a1j · · · a1n               +  k1 · · · kj · · · kn               an1 · · · anj · · · ann

           

This result can be seen by examining the terms of (8.1) for each of these determinants: A =


−1sgn
p a1p
1 a2p
2 · · ·
kp
k + kp
k  · · · anp
n p

 =
−1sgn
p a1p
1 a2p
2 · · · kp
k · · · anp
n p

+




−1sgn
p a1p
1 a2p
2 · · · + kp
k · · · anp
n = A1  + A2  

p

As a corollary to this, adding a scalar multiple of one row to another of a matrix does not change the value of the determinant.

306

CHAPTER 8

Determinants

COROLLARY 8.3

Let B be formed from A by adding ! times row i to row k. Then B = A. This result follows immediately from the preceding theorem by noting that row k of B is !aij + akj . Then   a11        a i1    B =     !ai1 + ak1        an1   a11 · · ·          ai1 · · ·    = !      ai1 · · ·          an1 · · ·

ann

                

a11  

···  

a1j  

···  

a1n  

ai1  

···  

aij  

···  

ain  

ak1  

···  

akj  

···  

akn  

an1

···

anj

···

ann

···  

a1j  

···  

a1n  

···  

aij  

···  

ain  

···  

!aij + akj  

···  

!ain + akn  

···

anj

a1j  

···  

a1n  

aij  

···  

ain  

aij  

···  

ain  

anj

···

ann

···                   +              

          = A        

In the last line, the first term is ! times a determinant with rows i and k identical, hence is zero. The second term is just A. We now know the effect of elementary row operations on a determinant. In summary: Type I operation—interchange of two rows. This changes the sign of the determinant. Type II operation—multiplication of a row by a scalar . This multiplies the determinant by  Type III operation—addition of a scalar multiple of one row to another row. This does not change the determinant. Recall that the transpose At of a matrix A is obtained by writing the rows of A as the columns of At . We claim that a matrix and its transpose have the same determinant.

THEOREM 8.6

A = At . For example, consider the 2 × 2 case:   a11 a12 A=  a21 a22

 At =

a11 a12

a21 a22

 

Then A = a11 a22 − a12 a21

and

At  = a11 a22 − a21 a12 = A 

8.4 Evaluation of Determinants by Elementary Row and Column Operations

307

A proof of this theorem consists of comparing terms of the determinants. If A = aij  then At = aji . Now, from the defining sum (8.1), A =




−1sgn
p a1p
1 a2p
2 · · · anp
n

p

and At  =




−1sgn
p
At ip
1
At 2p
2 · · ·
At np
n

p

=




−1sgn
q aq
11 aq
22 · · · aq
nn 

q

One can show that each term
−1sgn
p a1p
1 a2p
2 · · · anp
n in the sum for A is equal to a corresponding term
−1sgn
q aq
11 aq
22 · · · aq
nn in the sum for At . The key is to realize that, because q is a permutation of 1     n, we can rearrange the terms in the latter product to write them in increasing order of the first (row) index. This induces a permutation on the second (column) index, and we can match this term up with a corresponding term in the sum for A. We will not elaborate the details of this argument. One consequence of this result is that we can perform not only elementary row operations on a matrix, but also the corresponding elementary column operations, and we know the effect of each operation on the determinant. In particular, from the column perspective: If two columns of A are identical, or if one column is a zero column, then A = 0. Interchange of two columns of A changes the sign of the determinant. Multiplication of a column by a scalar  multiplies the determinant by  And addition of a scalar multiple of one column to another column does not change the determinant. These operations on rows and columns of a matrix, and their effect on the determinant of the newly formed matrix, form the basis for strategies to evaluate determinants.

SECTION 8.3

PROBLEMS

1. Let A = aij  be an n×n matrix and let  be any scalar. Let B = aij . Thus B is formed by multiplying each element of A by . Prove that B = n A. 2. Let A = aij  be an n × n matrix. Let  be a nonzero number. Form a new matrix B = i−j aij . How are A and

8.4

B related? Hint: It is useful to examine the 2 × 2 and 3 × 3 cases to get some idea of what B looks like. 3. An n × n matrix is skew-symmetric if A = −At . Prove that the determinant of a skew-symmetric matrix of odd order is zero.

Evaluation of Determinants by Elementary Row and Column Operations The use of elementary row and column operations to evaluate a determinant is predicated upon the following observation. If a row or column of an n × n matrix A has all zero elements except possibly for aij in row i and column j, then the determinant of A is
−1i+j aij times the determinant of the
n − 1 ×
n − 1 matrix obtained by deleting row i and column j from A.

308

CHAPTER 8

Determinants

This reduces the problem of evaluating an n × n determinant to one of evaluating a smaller determinant, having one less row and one less calumn. Here is a statement of this result, with (1) the row version and (2) the column version. THEOREM 8.7

1. Row Version   a11         ai−11   0  a  i+11        a

···  

a1j−1  

a1j  

a1j+1  

···  

a1n  

              

· · · ai−1j−1 ai−1j ai−1j+1 · · · ai−1n ··· 0 aij 0 ··· 0 · · · ai+1j−1 ai+1j ai+1j+1 · · · ai+1n             · · · a a a · · · a n1 nj−1 nj nj+1 nn   a11 · · · a1j−1 a1j+1 · · · a1n                  · · · a a · · · a a i−1j−1 i−1j+1 i−1n =
−1i+j aij  i−11  ai+11 · · · ai+1j−1 ai+1j+1 · · · ai+1n                 an1 · · · anj−1 anj+1 · · · ann

2. Column Version   a11         ai−11   ai1  a  i+11        a

···  

a1j−1  

0  

a1j+1  

· · · ai−1j−1 0 ai−1j+1 · · · aij−1 aij aij+1 · · · ai+1j−1 0 ai+1j+1         · · · a 0 a n1 nj−1 nj+1   a11 · · · a1j−1            ai−11 · · · ai−1j−1 i+j =
−1 aij   ai+11 · · · ai+1j−1           an1 · · · anj−1

···  

a1n  

··· ··· ···  

ai−1n ain ai+1n  

···

ann

             

              

a1j+1  

···  

a1n  

ai−1j+1 ai+1j+1  

··· ···  

ai−1n ai+1n  

anj+1

···

ann

             

This result suggests one strategy for evaluating a determinant. Given an n × n matrix A, use the row and/or column operations to obtain a new matrix B having at most one nonzero element in some row or column. Then A is a scalar multiple of B, and B is a scalar multiple of the
n − 1 ×
n − 1 determinant formed by deleting from B the row and column containing this nonzero element. We can then repeat this strategy on this
n − 1 ×
n − 1 matrix, eventually reducing the problem to one of evaluating a “small” determinant. Here is an illustration of this process.

8.4 Evaluation of Determinants by Elementary Row and Column Operations

309

EXAMPLE 8.1

Let

⎛ ⎜ ⎜ A=⎜ ⎜ ⎝

−6 −1 8 0 1

⎞ 0 1 3 2 5 0 1 7 ⎟ ⎟ 3 2 1 7 ⎟ ⎟ 1 5 −3 2 ⎠ 15 −3 9 4

We want to evaluate A. There are many ways to proceed with the strategy we are illustrating. To begin, we can exploit the fact that a13 = 1 and use elementary row operations to get zeros in the rest of column 3. Of course a23 = 0 to begin with, so we need only worry about column 3 entries in rows 3 4 5. Add
−2(row 1) to row 3, −5(row 1) to row 4 and 3(row 1) to row 5 to get ⎞ ⎛ −6 0 1 3 2 ⎜ −1 5 0 1 7 ⎟ ⎟ ⎜ ⎜ 3 0 −5 3 ⎟ B = ⎜ 20 ⎟ ⎝ 30 1 0 −18 −8 ⎠ −17 15 0 18 10 Because we have used Type III row operations, A = B  Further, by Theorem 8.7, B =
−11+3 b13 C =
1 C = C  where C is the 4 × 4 matrix obtained by deleting row 1 and column 3 of B: ⎛ ⎞ −1 5 1 7 ⎜ 20 3 −5 3 ⎟ ⎟ C=⎜ ⎝ 30 1 −18 −8 ⎠ −17 15 18 10 This is a 4 × 4 matrix, “smaller” than A. We will now apply the strategy to C. We can, for example, exploit the −1 entry in the 1 1 position of C, this time using column operations to get zeros in row 1, columns 2 3 4 of the new matrix. Specifically, add 5(column 1 ) to column 2, add column 1 to column 3, and add 7(column 1) to column 4 of C to get ⎛ ⎞ −1 0 0 0 ⎜ 20 103 15 143 ⎟ ⎟ D=⎜ ⎝ 30 151 12 202 ⎠ −17 −70 1 −109 Again, because we used Type III operations (this time on columns) of C, then C = D  But by the theorem, because we are using the element d11 = −1 as the single nonzero element of row 1, we have D =
−11+1 d11 E = − E 

CHAPTER 8

310

Determinants

in which E is the 3 × 3 matrix obtained by deleting row 1 and column 1 from D: ⎛ ⎞ 103 15 143 202 ⎠  E = ⎝ 151 12 −70 1 −109 To evaluate E, we can exploit the entry 3 2 −12(row 3) to row 2 to get ⎛ 1153 F = ⎝ 991 −70

entry e31 = 1. Add −15(row 3) to row 1 and ⎞ 0 1778 0 1510 ⎠  1 −109

Then E = F  By the theorem, using the only nonzero element f32 = 1 of column 2 of F, we have F =
−13+2
1 G = − G in which G is the 2 × 2 matrix obtained by deleting row 3 and column 2 of F:   1153 1778 G=  991 1510 At the 2 × 2 state, we evaluate the determinant directly: G =
1153
1510 −
1778
991 = −20 968 Working back, we now have A = B = C = D = − E = − F = G = −20 968 The method is actually quicker to apply than might appear from this example, because we included comments as we proceeded with the calculations.

SECTION 8.4

PROBLEMS

In each of Problems 1 through 10, use the strategy of this section to evaluate the determinant of the matrix. ⎛

−2 1. ⎝ 1 7 ⎛ 2 2. ⎝ 14 −13 ⎛ −4 3. ⎝ −2 2

4 6 0

⎞ 1 3 ⎠ 4

−3 1 −1 5 3 −2

⎞ 7 1 ⎠ 5 ⎞ 6 5 ⎠ 6

⎛

2 4. ⎝ 4 13 ⎛

17 5. ⎝ 1 14 ⎛

−3 ⎜ 1 ⎜ 6. ⎝ 7 2

−5 3 0

⎞ 8 8 ⎠ −4

−2 12 7

⎞ 5 0 ⎠, −7

3 −2 1 1

9 15 1 −1

⎞ 6 6 ⎟ ⎟ 5 ⎠ 3

8.5 Cofactor Expansions ⎛

0 ⎜ 6 ⎜ 7. ⎝ 1 4 ⎛

2 ⎜ 3 8. ⎜ ⎝ −2 4

8.5

1 −3 −5 8 7 1 0 8

1 2 1 2 −1 1 3 −1

⎛

⎞ −4 2 ⎟ ⎟ −2 ⎠ 2 ⎞ 0 8 ⎟ ⎟ 1 ⎠ 0

10 ⎜ 0 ⎜ 9. ⎝ 0 −2 ⎛ −7 ⎜ 1 10. ⎜ ⎝ 0 6

1 3 1 6 16 0 3 1

311

⎞ 2 9 ⎟ ⎟ 7 ⎠ 8

−6 3 1 8 2 0 −4 1

⎞ 4 5 ⎟ ⎟ 4 ⎠ −5

Cofactor Expansions Theorem 8.5 suggests the following. If we select any row i of a square matrix A, we can write     a11 a12 · · · · · · a1n   a11 a12                          ai1 ai2 · · · · · · ain  =  ai1 0                          an1 an2 · · · · · · ann   an1 an2     a11  a11 a12 · · · · · · a1n                          +  0 ai2 · · · · · · 0  + · · · +  0                      an1  an1 an2 · · · · · · ann 

··· ···     ··· ···     ··· ···

      0      ann  a1n  

a1n  

0  

··· ···     ··· ···    

an2

··· ···

ann

a12  

ain  

           

(8.2)

Each of the determinants on the right of equation (8.2) has a row in which every element but possibly one is zero, so Theorem 8.7 applies to each of these determinants. The first determinant on the right is
−1i+1 ai1 times the determinant of the matrix obtained by deleting row i and column 1 from A. The second determinant on the right is
−1i+2 ai2 times the determinant of the matrix obtained by deleting row i and column 2 from A. And so on, until the last matrix on the right is
−1i+n ain times the determinant of the matrix obtained by deleting row i and column n from A. We can put all of this more succinctly by introducing the following standard terminology.

DEFINITION 8.2

Minor

If A is an n × n matrix, the minor of aij is denoted Mij , and is the determinant of the
n − 1 ×
n − 1 matrix obtained by deleting row i and column j of A. Cofactor The number
−1i+j Mij is called the cofactor of aij .

We can now state the following formula for a determinant.

312

CHAPTER 8 THEOREM 8.8

Determinants Cofactor Expansion by a Row

If A is n × n, then for any integer i with i ≤ i ≤ n, A =

n 
−1i+j aij Mij 

(8.3)

j=1

This is just equation (8.2) in the notation of cofactors. The sum (8.3) is called the cofactor expansion of A by row i because it is the sum, across this row, of each matrix element times its cofactor. This yields A no matter which row is used. Of course, if some aik = 0 then we need not calculate that term in equation (8.3), so it is to our advantage to expand by a row having as many zero elements as possible. The strategy of the preceding subsection was to create such a row using row and column operations, resulting in what was a cofactor expansion, by a row having only one (possibly) nonzero element.

EXAMPLE 8.2

Let ⎛

−6 3 A = ⎝ 12 −5 2 4

⎞ 7 −9 ⎠  −6

If we expand by row 1, we get A =

3 
−11+j M1j j=1

    −5 −9    +
−11+2
3  12 =
−11+1
−6    2 4 −6    12 −5   +
−11+3
7  2 4 

 −9  −6 

=
−6
30 + 36 − 3
−72 + 18 + 7
−48 + 10 = 172 Just for illustration, expand by row 3:  3   3+j 3+1  A =
−1 a3j M3j =
−1 2 

 7  −9  j=1     −6  −6 7  3+3  +
−1 +
−13+2
4 
−6   12 12 −9 3 −5

 3  −5 

= 2
−27 + 35 − 4
54 − 84 − 6
30 − 36 = 172

Because, for purposes of evaluating determinants, row and column operations can both be used, we can also develop a cofactor expansion of A by column j. In this expansion, we move down a column of a matrix and sum each term of the column times its cofactor.

8.5 Cofactor Expansions THEOREM 8.9

313

Cofactor Expansion by a Column

Let A be an n × n matrix. Then for any j with 1 ≤ j ≤ n, A =

n 


−1i+j aij Mij 

(8.4)

i=1

This differs from the expansion (8.3) in that the latter expands across a row, while the sum (8.4) expands down a column. All of these expansions, by any row or column of A, yield A.

EXAMPLE 8.3

Consider again

⎛

−6 A = ⎝ 12 2

3 −5 4

⎞ 7 −9 ⎠  −6

Expanding by column 1 gives us A =

  3  −5 −9   
−1i+1 ai1 Mi1 =
−11+1
−6  4 −6  i=1     3  3 7  7 2+1 3+1  +
−1
12  +
−1
2  4 −6  −5 −9

   

=
−6
30 + 36 − 12
−18 − 28 + 2
−27 + 35 = 172 If we expand by column 2 we get

   12 −9    A =
−1 ai2 Mi2 =
−1
3  2 −6  i=1     −6  −6 7  3+2  +
−12+2
−5  +
−1
4   12 2 −6 3 

i+2

1+2

=
−3
−72 + 18 − 5
36 − 14 − 4
54 − 84 = 172

SECTION 8.5

PROBLEMS

In Problems 1–10, use cofactor expansions, combined with elementary row and column operations when this is useful, to evaluate the determinant of the matrix. ⎞ ⎛ −4 2 −8 1 0 ⎠ 1. ⎝ 1 1 −3 0 ⎞ ⎛ 1 1 6 2. ⎝ 2 −2 1 ⎠ 3 −1 4

⎛

7 3. ⎝ 1 −3 ⎛ 5 4. ⎝ −1 −2 ⎛ −5 ⎜ 2 ⎜ 5. ⎝ 4 1

−4 1 −2

⎞ 1 4 ⎠ 0 ⎞ 3 6 ⎠ 4

0 −1 4 −1

1 3 −5 6

−3 −2 1

⎞ 6 7 ⎟ ⎟ −8 ⎠ 2

 7  −9 

CHAPTER 8

314

Determinants

⎞ 4 3 −5 6 ⎜ 1 −5 15 2 ⎟ ⎟ 6. ⎜ ⎝ 0 −5 1 7 ⎠ 8 9 0 15 ⎞ ⎛ −3 1 14 1 16 ⎠ 7. ⎝ 0 2 −3 4 ⎞ ⎛ 14 13 −2 5 ⎜ 7 1 1 7 ⎟ ⎟ 8. ⎜ ⎝ 0 2 12 3 ⎠ 1 −6 5 2 ⎞ ⎛ −5 4 1 7 ⎜ −9 3 2 −5 ⎟ ⎟ 9. ⎜ ⎝ −2 0 −1 1 ⎠ 1 14 0 3 ⎛ −8 5 1 7 2 ⎜ 0 1 3 5 −6 ⎜ 1 5 3 10. ⎜ ⎜ 2 2 ⎝ 0 4 3 7 2 1 1 −7 −6 5 ⎛

11. Show that   1    1    1 !

8.6

2 2 !2

This is Vandermonde’s determinant. This and the next problem are best done with a little thought in using facts about determinants, rather than a brute-force approach. 12. Show         !   

that  !  

!   

     =
 +  + ! + 
 −  +  − !       0 1 −1 1     1 !    ×    1      1   ! 

   !

13. Let A be a square matrix such that A−1 = At . Prove that A = ±1.

⎞ ⎟ ⎟ ⎟ ⎟ ⎠

14. Prove that three points
x1  y1 ,
x2  y2  and
x3  y3  are collinear (on the same straight line) if and only if    1 x1 y1     1 x2 y2  = 0    1 x3 y3 

    =
 − 
! − 
 − !  

Hint: This determinant is zero exactly when one row or column is a linear combination of the other two.

Determinants of Triangular Matrices The main diagonal of a square matrix A consists of the elements a11 , a22 ,    , ann . We call A upper triangular if all the elements below the main diagonal are zero. That is, aij = 0 if i > j. Such a matrix has the appearance ⎛

a11 0 0  

⎜ ⎜ ⎜ ⎜ A=⎜ ⎜ ⎜ ⎝ 0 0

a12 a22 0  

a13 a23 a33  

··· ··· ···  

a1n−1 a2n−1 a3n−1  

a1n a2n a3n  

0 0

0 0

··· 0

an−1n−1 0

an−1n ann

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

If we expand A by cofactors down the first column, we have      A = a11    

a22 0  

a23 a33  

··· ···  

a2n−1 a3n−1  

a2n a3n  

0 0

0 0

··· ···

an−1n−1 0

an−1n ann

         

8.7 A Determinant Formula for a Matrix Inverse

315

and the determinant on the right is again upper triangular, so expand by its first column to get    a33 a34 · · · a3n     0 a44 · · · a3n    A = a11 a22         ···     0 0 · · · ann  with another upper triangular determinant on the right. Continuing in this way, we obtain A = a11 a22 · · · ann  The determinant of an upper triangular matrix is the product of its main diagonal elements. The same conclusion holds for lower triangular matrices (all elements above the main diagonal are zero). Now we can expand the determinant along the top row, each time obtaining just one minor that is again lower triangular.

EXAMPLE 8.4

         

15 0 0 0 0

Evaluate the following determinants.   0 0 0 0 0   −4   7 0 0 0 0   12   2 0 0 0   3 −4 1.   1 1 −2 0 0   0   1 5 0   1 −4 16   10 −4 16 1 17 4

8.7

6 0 0 0 0 0

1 −4 0 0 0 0

−1 2 −5 0 0 0

2 2 10 14 0 0

4 −6 √ 2 0 0

 7 3  3 9  √ √ 15 −4  =
15
12 2e = 180 2e  22  0 e 

PROBLEMS

SECTION 8.6

      2.     

−7 12 0 0 0

2 −3 1 0 13 0

1 1 −7 0 −4 3

  3   2  3.  17   22  43

0 −6 14 −2 12

0 0 2 15 1

0 0 0 8 −1

0 0 0 0 5

        

          

A Determinant Formula for a Matrix Inverse Determinants can be used to tell whether a matrix is singular or nonsingular. In the latter case, there is a way of writing the inverse of a matrix by using determinants.

316

CHAPTER 8

Determinants

First, here is a simple test for nonsingularity. We will use the fact that we reduce a matrix by using elementary row operations, whose effects on determinants are known (Type I operations change the sign, Type II operations multiply the determinant by a nonzero constant, and Type III operations do not change the determinant at all). This means that, for any square matrix A, A =  AR  for some nonzero constant . THEOREM 8.10

Let A be an n × n matrix. Then A is nonsingular if and only if A = 0. Suppose first that A = 0. Since A =  AR  for some nonzero constant , AR can have no zero row, so AR = In . Then rank
A = n, so A is nonsingular by Theorems 7.26 and 7.15. Conversely, suppose A is nonsingular. Then AR = In . Then A =  AR  = 0.

Proof

Using this result, we can give a short proof of Theorem 7.25(5). Suppose A and B are n × n matrices, and AB is singular. Then AB = A B = 0 so A = 0 or B = 0, hence either A or B (or possibly both) must be singular. We will now write a formula for the inverse of a square matrix, in terms of cofactors of the matrix. THEOREM 8.11

Let A be an n × n nonsingular matrix. Define an n × n matrix B by putting bij =

1
−1i+j Mji  A

Then, B = A−1 . That is, the i j element of A−1 is the cofactor of aji (not aij ), divided by the determinant of A. Proof

By the way B is defined, the i j element of AB is
ABij =

n 

aik bkj =

k=1

n 1 
−1j+k aik Mjk  A k=1

Now examine the sum on the right. If i = j, we get
ABii =

n 1 
−1i+k aik Mik A k=1

and the summation is exactly the cofactor expansion of A by row i. Therefore
ABii =

A = 1 A

If i = j, then the summation in the expression for
ABij is the cofactor expansion, by row j, of the determinant of the matrix formed from A by replacing row j by row i. But this matrix then has two identical rows, hence has determinant zero. Then
ABij = 0 if i = j, and we conclude that AB = In . A similar argument shows that BA = In , hence B = A−1 .

8.7 A Determinant Formula for a Matrix Inverse

317

This method of computing a matrix inverse is not as efficient in general as the reduction method discussed previously. Nevertheless, it works well for small matrices, and in some discussions it is useful to have a formula for the elements of a matrix inverse.

EXAMPLE 8.5

Let

⎛

−2 A=⎝ 6 2 Then

  −2   6   2

⎞ 4 1 3 −3 ⎠  9 −5

4 1 3 −3 9 −5

    = 120  

so A is nonsingular. Compute the nine elements of the inverse matrix B:   1 1 1  3 −3  12 b11 = M11 = =  =   9 −5 120 120 120 10 b12 =

 1 1  4
−1M21 = − 120 120  9

 1 1  4 b13 = M = 120 31 120  3

 1 1  =−  −3  8

 1 1  6 b21 = − M =− 120 12 120  2 b22 =

 1  −2 120  2

 1 1  6 M13 = 120 120  2

Then

 1  = 0 −3 

 3  2 =  9  5

 1 1  −2 b32 = − M =− 120 23 120  2  1  −2 b33 = 120  6

 −3  1 =  −5  5

 1 1  =  −5  15

 1 1  −2 b23 = − M =− 120 32 120  6 b31 =

 29 1   = −5  120

 4  13 =  9  60

 1 4  =−   3 4 ⎛

⎜ B = A−1 = ⎝

1 10 1 5 2 5

29 120 1 15 13 60

−1 8

⎞

⎟ 0 ⎠

−1 4

CHAPTER 8

318

SECTION 8.7

Determinants

PROBLEMS

In each of Problems 1 through 10, use Theorem 8.10 to determine whether the matrix is nonsingular. If it is, use Theorem 8.11 to find its inverse.   2 −1 1. 1 6   3 0 2. 1 4   −1 1 3. 1 4   2 5 4. −7 −3 ⎞ ⎛ 6 −1 3 1 −4 ⎠ 5. ⎝ 0 2 2 −3 ⎞ ⎛ −14 1 −3 2 −1 3 ⎠ 6. ⎝ 1 1 7

8.8

⎛

0 7. ⎝ 2 1 ⎛

11 8. ⎝ 0 4 ⎛

3 ⎜ 4 9. ⎜ ⎝ −2 13 ⎛ 7 ⎜ 8 10. ⎜ ⎝ 1 3

−4 −1 −1 0 1 −7

⎞ 3 6 ⎠ 7 ⎞ −5 0 ⎠ 9

1 6 1 0

−2 −3 7 1

⎞ 1 9 ⎟ ⎟ 4 ⎠ 5

−3 2 5 −2

−4 0 −1 −5

⎞ 1 0 ⎟ ⎟ 7 ⎠ 9

Cramer’s Rule Cramer’s rule is a determinant formula for solving a system of equations AX = B when A is n × n and nonsingular. In this case, the system has the unique solution X = A−1 B. We can, therefore, find X by computing A−1 and then A−1 B. Here is another way to find X.

THEOREM 8.12

Cramer’s Rule

Let be a nonsingular n × n matrix of numbers. Then the unique solution of AX = B is ⎛ A⎞ x1 ⎜ x2 ⎟ ⎜ ⎟ ⎜  ⎟, where ⎝  ⎠ xn xk =

1 A
k B A

and A
k B is the matrix obtained from A by replacing column k of A by B. Here is a heuristic argument to suggest why this works. Let ⎛ ⎞ b1 ⎜ b2 ⎟ ⎜ ⎟ B = ⎜  ⎟ ⎝  ⎠ bn

8.8 Cramer’s Rule

319

Multiply column k of A by xk . The determinant of the resulting matrix is xk A, so    a11 a12 · · · a1k xk · · · a1n     a21 a22 · · · a2k xk · · · a2n    xk A =                     an1 an2 · · · ank xk · · · ann  For each j = k, add xj times column j to column k. This Type III operation does not change the value of the determinant, and we get    a11 a12 · · · a11 x1 + a12 x2 + · · · + a1n xn · · · a1n     a21 a22 · · · a21 x1 + a22 x2 + · · · + a2n an · · · a2n    xk A =                   an1 an2 · · · an1 x1 + an2 x2 + · · · + ann xn · · · ann     a11 a12 · · · b1 · · · a1n     a21 a22 · · · b2 · · · a2n    =        = A
k B             an1 an2 · · · bn · · · ann  Solving for xk yields the conclusion of Cramer’s Rule.

EXAMPLE 8.6

Solve the system x1 − 3x2 − 4x3 = 1 −x1 + x2 − 3x3 = 14 x2 − 3x3 = 5 The matrix of coefficients is

⎛

⎞ 1 −3 −4 1 −3 ⎠  A = ⎝ −1 0 1 −3

We find that A = 13. By Cramer’s rule,   1 −3 1  14 1 x1 = 13  5 1   1 1 1  −1 14 x2 = 13  0 5   1 −3 1  −1 1 x3 = 13  0 1

    = − 117 = −9  13   −4  10 −3  = −  13 −3   1  25 14  = −  13 5  −4 −3 −3

CHAPTER 8

320

Determinants

Cramer’s rule is not as efficient as the Gauss-Jordan reduction. Gauss-Jordan also applies to homogeneous systems, and to systems with different numbers of equations than unknowns. However, Cramer’s rule does provide a formula for the solution, and this is useful in some contexts.

SECTION 8.8

PROBLEMS

In each of Problems 1 through 10, either find the solution by Cramer’s rule or show that the rule does not apply. 1.

6.

6x1 + 4x2 − x3 + 3x4 − x5 = 7 x1 − 4x2 + x5 = −5 x1 − 3x2 + x3 − 4x5 = 0 −2x1 + x3 − 2x5 = 4 x3 − x4 − x5 = 8

7.

2x1 − 4x2 + x3 − x4 = 6 x2 − 3x3 = 10 x1 − 4x3 = 0 x2 − x3 + 2x4 = 4

8.

2x1 − 3x2 + x4 = 2 x2 − x3 + x4 = 2 x3 − 2x4 = 5 x1 − 3x2 + 4x3 = 0

9.

14x1 − 3x3 = 5 2x1 − 4x3 + x4 = 2 x1 − x2 + x3 − 3x4 = 1 x3 − 4x4 = −5

10.

x2 − 4x4 = 18 x1 − x2 + 3x3 = −1 x1 + x2 − 3x3 + x4 = 5 x2 + 3x4 = 0

15x1 − 4x2 = 5 8x1 + x2 = −4

2.

x1 + 4x2 = 3 x1 + x2 = 0

3.

8x1 − 4x2 + 3x3 = 0 x1 + 5x2 − x3 = −5 −2x1 + 6x2 + x3 = −4

4.

5x1 − 6x2 + x3 = 4 −x1 + 3x2 − 4x3 = 5 2x1 + 3x2 + x3 = −8

5.

x1 + x2 − 3x3 = 0 x2 − 4x3 = 0 x1 − x2 − x3 = 5

8.9

The Matrix Tree Theorem In 1847, G.R. Kirchhoff published a classic paper in which he derived many of the electrical circuit laws that bear his name. One of these is the matrix tree theorem, which we will discuss now. Figure 8.1 shows a typical electrical circuit. The underlying geometry of the circuit is shown in Figure 8.2. This diagram of points and connecting lines is called a graph, and was seen in the context of the movement of atoms in crystals in Section 7.1.5. A labeled graph has symbols attached to the points. Some of Kirchhoff’s results depend on geometric properties of the circuit’s underlying graph. One such property is the arrangement of the closed loops. Another is the number of spanning trees in the labeled graph. A spanning tree is a collection of lines in the graph forming no closed loops, but containing a path between any two points of the graph. Figure 8.3 shows a labeled graph and two spanning trees in this graph. Kirchhoff derived a relationship between deteminants and the number of labeled trees in a graph.

8.9 The Matrix Tree Theorem

321

V

FIGURE 8.1

v2

v1

v3

FIGURE 8.2

v1

v2

v7 v10

v8

v5

v7

v10

v8

v10

v9 v4

v6

v8 v9

v5

v6

v4

v5

FIGURE 8.3 A labeled graph and two of its spanning trees.

THEOREM 8.13

Matrix Tree Theorem

Let G be a graph with vertices labeled v1      vn . Form an n × n matrix T = tij  as follows. If i = j, then tii is the number of lines to vi in the graph. If i = j, then tij = 0 if there is no line between vi and vj in G, and tij = −1 if there is such a line. Then, all cofactors of T are equal, and their common value is the number of spanning trees in G.

EXAMPLE 8.7

For the labeled graph of Figure 8.4, T is ⎛ 3 −1 ⎜ −1 3 ⎜ ⎜ 0 −1 ⎜ T =⎜ ⎜ 0 −1 ⎜ 0 0 ⎜ ⎝ −1 0 −1 0 G v7

the 7 × 7 matrix 0 0 0 −1 −1 −1 −1 0 0 0 3 −1 0 −1 0 −1 4 −1 0 −1 0 −1 3 −1 −1 −1 0 −1 4 −1 0 −1 −1 −1 4 v1

v2 v3

v6 v5

v3

v2

v1

v7

v9 v6

v3

v4

FIGURE 8.4 Graph G.

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

v4

322

CHAPTER 8

Determinants

Evaluate any cofactor of T. For example, covering up   3 −1 −1  3 −1  −1  4  −1 −1 1+1
−1 M11 =  0 −1  0  0 −1 0   0 0 −1

row 1 and column 1, we have  0 0 0   0 −1 0   −1 0 −1   3 −1 −1  −1 4 −1  −1 −1 4 

= 386 Evaluation of any cofactor of T yields the same result. Even with this small graph, it would clearly be impractical to enumerate the spanning trees by attempting to list them all.

PROBLEMS

SECTION 8.9

1. Find the number of spanning trees in the graph of Figure 8.5.

1

4. Find the number of spanning trees in the graph of Figure 8.8.

5 1

2 4

6

3

2. Find the number of spanning trees in the graph of Figure 8.6.

1 6

3

4

FIGURE 8.8

5. Find the number of spanning trees in the graph of Figure 8.9.

5

4

3 5

FIGURE 8.5

2

2

2 3

1

FIGURE 8.6

3. Find the number of spanning trees in the graph of Figure 8.7.

1 6

4 5

FIGURE 8.9

2 3

5

6

4

FIGURE 8.7

6. A complete graph on n points consists of n points with a line between each pair of points. This graph is often denoted Kn . With the points labeled 1 2     n, show that the number of spanning trees in Kn is nn−2 for n = 3 4    .

CHAPTER

9

EIGENVALUES AND EIGENVECTORS DIAGONALIZATION OF MATRICES ORTHOGONAL AND SYMMETRIC MATRICES QUADRATIC FORMS UNITARY HERMITIAN, AND SKEW HERMITIAN MATRICES EIGENVALUES

Eigenvalues, Diagonalization, and Special Matrices

Suppose A is an n × n matrix of real numbers. If we write an n-vector E as a column ⎛ ⎜ ⎜ E=⎜ ⎝

1 2  

⎞ ⎟ ⎟ ⎟ ⎠

n then AE is an n × 1 matrix which we may also think of as an n-vector. We may therefore consider A as an operator that moves vectors about in Rn . Because A
aE1 + bE2  = aAE1 + bAE2 , A is called a linear operator. Vectors have directions associated with them. Depending on A, the direction of AE will generally be different from that of E. It may happen, however, that for some vector E, AE and E are parallel. In this event there is a number  such that AE = E. Then  is called an eigenvalue of A, with E an associated eigenvector. The idea of an operator moving a vector to a parallel position is simple and geometrically appealing. It also has powerful ramifications in a variety of contexts. Eigenvalues contain important information about the solutions of systems of differential equations, and in models of physical phenomena may they have physical significance as well (such as the modes of vibration of a mechanical system, or the energy states of an atom). 323

324

9.1

CHAPTER 9

Eigenvalues, Diagonalization, and Special Matrices

Eigenvalues and Eigenvectors Let A be an n × n matrix of real or complex numbers.

DEFINITION 9.1

Eigenvalues and Eigenvectors

A real or complex number  is an eigenvalue of A if there is a nonzero n × 1 matrix (vector) E such that AE = E Any nonzero vector E satisfying this relationship is called an eigenvector associated with the eigenvalue .

Eigenvalues are also known as characteristic values of a matrix, and eigenvectors can be called characteristic vectors. We will typically write eigenvectors as column matrices and think of them as vectors in Rn . If an eigenvector has complex components, we may think of it as a vector in C n , which consists of n-tuples of complex numbers. Since an eigenvector must be nonzero vector, at least one component is nonzero. If  is a nonzero scalar and AE = E, then A
E = 
AE = 
E = 
E This means that nonzero scalar multiples of eigenvectors are again eigenvectors.

EXAMPLE 9.1

Since



1 0

0 0



0 4 



 =

0 0



 =0

0 4

 

 0 an associated eigenvector. Although the zero 4 vector cannot be an  eigenvector, the number zero can be an eigenvalue of a matrix. For any 0 scalar  = 0, is also an eigenvector associated with the eigenvalue 0. 4

0 is an eigenvalue of this matrix, with

EXAMPLE 9.2

Let

⎛

1 A=⎝ 0 0

⎞ −1 0 1 1 ⎠ 0 −1

9.1 Eigenvalues and Eigenvectors

325

⎛

⎞ 6 Then 1 is an eigenvalue with associated eigenvector ⎝ 0 ⎠, because 0 ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 6 6 6 A⎝ 0 ⎠ = ⎝ 0 ⎠ = 1⎝ 0 ⎠ 0 0 0 ⎛

⎞  Because any nonzero multiple of an eigenvector is an eigenvector, then ⎝ 0 ⎠ is also an 0 eigenvector associated with eigenvalue 1, for any nonzero number⎛. ⎞ 1 Another eigenvalue of A is −1, with associated eigenvector ⎝ 2 ⎠, because −4 ⎞ ⎞⎛ ⎛ 1 1 −1 0 ⎝ 0 1 1 ⎠⎝ 2 ⎠ −4 0 0 −1 ⎛ ⎞ ⎛ ⎞ −1 1 = ⎝ −2 ⎠ = −1 ⎝ 2 ⎠  4 −4 ⎛ ⎞  Again, any vector ⎝ 2 ⎠, with  = 0, is an eigenvector associated with −1. −4 We would like a way of finding all of the eigenvalues of a matrix A. The machinery to do this is at our disposal, and we reason as follows. For  to be an eigenvalue of A, there must be an associated eigenvector E, and AE = E. Then E − AE = O, or In E − AE = O The identity matrix was inserted so we could write the last equation as
In − AE = O This makes E a nontrivial solution of the n × n system of linear equations
In − AX = O But this system can have a nontrivial solution if and only if the coefficient matrix has determinant zero, that is, In − A = 0. Thus  is an eigenvalue of A exactly when In − A = 0. This is the equation     − a11 −a12 · · · −a1n     −a21  − a22 · · · −a2n      = 0              −an1 −an2 · · ·  − ann  When the determinant on the left is expanded, it is a polynomial of degree n in , called the characteristic polynomial of A. The roots of this polynomial are the eigenvalues of A. Corresponding to any root , any nontrivial solution E of
In − AX = O is an eigenvector associated with . We will summarize these conclusions.

326

CHAPTER 9

Eigenvalues, Diagonalization, and Special Matrices

THEOREM 9.1

Let A be an n × n matrix of real or complex numbers. Then 1.  is an eigenvalue of A if and only if In − A = 0. 2. If  is an eigenvalue of A, then any nontrivial solution of
In − AX = O is an associated eigenvector.

DEFINITION 9.2

Characteristic Polynomial

The polynomial In − A is the characteristic polynomial of A, and is denoted pA
.

If A is n × n, then pA
 is an nth degree polynomial with real or complex coefficients determined by the elements of A. This polynomial therefore has n roots, though some may be repeated. An n × n matrix A always has n eigenvalues 1      n , in which each eigenvalue is listed according to its multiplicity as a root of the characteristic polynomial. For example, if pA
 =
 − 1
 − 32
 − i4 we list 7 eigenvalues: 1 3 3 i i i i. The eigenvalue 3 has multiplicity 2 and i has multiplicity 4.

EXAMPLE 9.3

Let

⎛

1 −1 1 A=⎝ 0 0 0

⎞ 0 1 ⎠ −1

as in Example 9.2. The characteristic polynomial is   −1 1  −1 pA
 =  0  0 0

0 −1 +1

     

=
 − 12
 + 1 The eigenvalues of A are 1 1 −1. To find eigenvectors associated with eigenvalue 1, solve ⎛ ⎞ 0 1 0
1I3 − AX = ⎝ 0 0 −1 ⎠ X = O 0 0 2 This has general solution

⎛

⎞  ⎝ 0 ⎠ 0

and these are the eigenvectors associated with eigenvalue 1, with  = 0

9.1 Eigenvalues and Eigenvectors For eigenvectors associated with −1, solve ⎛ −2 1
−1I3 − AX = ⎝ 0 −2 0 0 The general solution is

327

⎞ 0 −1 ⎠ X = O 0

⎞  ⎝ 2 ⎠ −4 ⎛

and these are the eigenvectors associated with eigenvalue −1, as long as  = 0

EXAMPLE 9.4

Let

 A=

1 −2 2 0

 

The characteristic polynomial is       1 0 1 −2   pA
x =  − 0 1 2 0     −1 2   = 
 − 1 + 4 = 2 −  + 4 =  −2   √ √ This has roots
1 + 15i/2 and
1 − 15i/2, and these are the eigenvalues of A. Even though A has real elements, the eigenvalues may√be complex. To find eigenvectors associated with
1 + 15i/2, solve the system
I2 − AX = O, which for this  is  √     1 + 15i 1 0 1 −2 − X = O 0 1 2 0 2 This is the system



√ 1+ 15i 2

−1

−2 or



2

√

1+ 15i 2

x1 x2



 =

0 0



√ −1 + 15i x1 + 2x2 = 0 2 √ 1 + 15i −2x1 + x2 = 0 2

We find the general solution of this system to be   1  1−√15i  4

and this is an eigenvector associated with the eigenvalue

√ 1+ 15i 2

for any nonzero scalar .

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Corresponding to the eigenvalue 1− 2 15i solve the system   1−√15i −1 2 2 X = O √ 1− 15i −2 2 obtaining the general solution

 

1

√ 1+ 15i 4

 

This is an eigenvector corresponding to the eigenvalue

√ 1− 15i 2

for any  = 0.

Finding the eigenvalues of a matrix is equivalent to finding the roots of an nth degree polynomial, and if n ≥ 3 this may be difficult. There are efficient computer routines which are usually based on the idea of putting the matrix through a sequence of transformations, the effect of which on the eigenvalues is known. This strategy was used previously to evaluate determinants. There are also approximation techniques, but these are sensitive to error. A number that is very close to an eigenvalue may not behave like an eigenvalue. We will conclude this section with a theorem due to Gerschgorin. If real eigenvalues are plotted on the real line, and complex eigenvalues as points in the plane, Gerschgorin’s theorem enables us to delineate regions of the plane containing the eigenvalues.

9.1.1 THEOREM 9.2

Gerschgorin’s Theorem Gerschgorin

Let A be an n × n matrix of real or complex numbers. For k = 1     n, let rk =

n    akj   j=1j=k

Let Ck be the circle of radius rk centered at
k  k , where akk = k + ik . Then each eigenvalue of A, when plotted as a point in the complex plane, lies on or within one of the circles C1      Cn . The circles Ck are called Gerschgorin circles. For the radius of Ck , read across row k and add the magnitudes of the row elements, omitting the diagonal element akk . The center of Ck is akk , plotted as a point in the complex plane. If the Gerschgorin circles are drawn and the disks they bound are shaded, then we have a picture of a region containing all of the eigenvalues of A.

EXAMPLE 9.5

Let

⎛

12i ⎜ 1 A=⎜ ⎝ 4 1 − 3i

1 9 −6 2 + i 1 −1 −9 1

⎞ −4 −1 ⎟ ⎟ 4i ⎠ 4 − 7i

A has characteristic polynomial pA
 = 4 +
3 − 5i3 +
18 − 4i2 +
290 + 90i + 1374 − 1120i

9.1 Eigenvalues and Eigenvectors

329

It is not clear what the roots of this polynomial are. Form the Gerschgorin circles. Their radii are: r1 = 1 + 9 + 4 = 14 √ √ r2 = 1 + 5 + 1 = 2 + 5 r3 = 4 + 1 + 4 = 9 and

√ 10 + 9 + 1 = 10 + 10 √ C1 has radius 14 and center
0 12, C2 has radius √ 2 + 5 and center
−6 0, C3 has radius 9 and center
−1 0 and C4 has radius 10 + 10 and center
4 −7. Figure 9.1 shows the Gerschgorin circles containing the eigenvalues of A. r4 =

√

y (0, 12)

(
6, 0) x (4,
7)

(
1, 0)

FIGURE 9.1 Gerschgorin circles.

Gerschgorin’s theorem is not intended as an approximation scheme, since the Gerschgorin circles may have large radii. For some problems, however, just knowing some information about possible locations of eigenvalues can be important. For example, in studies of the stability of fluid flow, it is important to know whether there are eigenvalues in the right half-plane.

SECTION 9.1

PROBLEMS

In each of Problems 1 through 16, (a) find the eigenvalues of the matrix, (b) corresponding to each eigenvalue, find an eigenvector, and (c) sketch the Gerschgorin circles and (approximately) locate the eigenvalues as points in the plane.  1.  2.  3.

1 2

3 1

0 4

−5 1

0 2

 5.  6.



−2 1

 4.



1 2 0 0

−2 4  −6 2  1 0



⎞ 2 0 0 7. ⎝ 1 0 2 ⎠ 0 0 3 ⎞ ⎛ −2 1 0 0 ⎠ 8. ⎝ 1 3 0 0 −1 ⎛



6 −3

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⎞ −3 1 1 9. ⎝ 0 0 0 ⎠ 0 1 0 ⎞ ⎛ 0 0 −1 1 ⎠ 10. ⎝ 0 0 2 0 0 ⎞ ⎛ −14 1 0 0 2 0 ⎠ 11. ⎝ 1 0 2 ⎞ ⎛ 3 0 0 12. ⎝ 1 −2 −8 ⎠ 0 −5 1 ⎞ ⎛ 1 −2 0 0 0 ⎠ 13. ⎝ 0 −5 0 7 ⎞ ⎛ −2 1 0 0 ⎜ 1 0 0 1 ⎟ ⎟ 14. ⎜ ⎝ 0 0 0 0 ⎠ 0 0 0 0 ⎞ ⎛ −4 1 0 1 ⎜ 0 1 0 0 ⎟ ⎟ 15. ⎜ ⎝ 0 0 2 0 ⎠ 1 0 0 3 ⎛

9.2

⎛

5 ⎜ 0 ⎜ 16. ⎝ 0 0

1 1 0 0

0 0 0 0

⎞ 9 9 ⎟ ⎟ 9 ⎠ 0



   !  and ! are real numbers, are real. ⎛   18. Show that the eigenvalues of ⎝   !  if all of the matrix elements are real. 17. Show that the eigenvalues of

 , in which , ⎞ !  ⎠ are real, "

19. Let  be an eigenvalue of A with eigenvector E. Show that, for any positive integer k, k is an eigenvalue of Ak , with eigenvector E. 20. Let  be an eigenvalue of A with eigenvector E, and  an eigenvalue of A with eigenvector L. Suppose  = . Show that E and L are linearly independent as vectors in Rn . 21. Let A be an n × n matrix. Prove that the constant term of pA
x is
−1n A. Use this to show that any singular matrix must have zero as one of its eigenvalues.

Diagonalization of Matrices We have referred to the elements aii of a square matrix as its main diagonal elements. All other elements are called off-diagonal elements.

DEFINITION 9.3

Diagonal Matrix

A square matrix having all off-diagonal elements equal to zero is called a diagonal matrix.

We often write a diagonal matrix having main diagonal elements d1      dn as ⎛ ⎜ ⎜ ⎜ ⎝

d1

O d2 

O



⎞ ⎟ ⎟ ⎟ ⎠

dn

with O in the upper right and lower left corners to indicate that all off-diagonal elements are zero. Here are some properties of diagonal matrices that make them pleasant to work with.

9.2 Diagonalization of Matrices

331

THEOREM 9.3

Let

⎛ ⎜ ⎜ D=⎜ ⎝

O

d1 d2  O



⎞

⎛

⎟ ⎟ ⎟ ⎠

⎜ ⎜ and W = ⎜ ⎝

dn

O

w1 w2 



O

⎞ ⎟ ⎟ ⎟ ⎠

wn

Then 1.

⎛ ⎜ ⎜ DW = WD = ⎜ ⎝

d 1 w1

O d2 w2 

O



⎞ ⎟ ⎟ ⎟ ⎠

dn wn

2. D = d1 d2 · · · dn . 3. D is nonsingular if and only if each main diagonal element is nonzero. 4. If each dj = 0, then ⎛ ⎞ 1/d1 O ⎜ ⎟ 1/d2 ⎜ ⎟ D−1 = ⎜ ⎟  ⎝ ⎠  1/dn

O 5. The eigenvalues of D are its main diagonal elements. 6. An eigenvector associated with dj is ⎞ ⎛ 0 ⎜  ⎟ ⎜  ⎟ ⎟ ⎜ ⎜ 0 ⎟ ⎟ ⎜ ⎜ 1 ⎟ ⎟ ⎜ ⎜ 0 ⎟ ⎟ ⎜ ⎜  ⎟ ⎝  ⎠ 0 with 1 in row j and all other elements zero.

We leave a proof of these conclusions to the student. Notice that (2) follows from the fact that a diagonal matrix is upper (and lower) triangular. “Most” square matrices are not diagonal matrices. However, some matrices are related to diagonal matrices in a way that enables us to utilize the nice features of diagonal matrices.

DEFINITION 9.4

Diagonalizable Matrix

An n × n matrix A is diagonalizable if there exists an n × n matrix P such that P−1 AP is a diagonal matrix. When such P exists, we say that P diagonalizes A.

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The following theorem not only tells us when a matrix is diagonalizable, but also how to find a matrix P that diagonalizes it. THEOREM 9.4

Diagonalizability

Let A be an n × n matrix. Then A is diagonalizable if it has n linearly independent eigenvectors. Further, if P is the n × n matrix having these eigenvectors as columns, then P−1 AP is the diagonal matrix having the corresponding eigenvalues down its main diagonal. Here is what this means. Suppose 1      n are the eigenvalues of A (some possibly repeated), and V1      Vn are corresponding eigenvectors. If these eigenvectors are linearly independent, we can form a nonsingular matrix P using Vj as column j. It is the linear independence of the eigenvectors that makes P nonsingular. We claim that P−1 AP is the diagonal matrix having the eigenvalues of A down its main diagonal, in the order corresponding to the order the eigenvectors were listed as columns of P.

EXAMPLE 9.6

Let



−1 0

A=



4 3

 

A has eigenvalues −1 3 and corresponding eigenvectors  P=

1 0

1 1

1 0



 and

1 1

 respectively. Form

 

Because the eigenvectors are linearly independent, this matrix is nonsingular (note that P = 0). We find that   1 −1 −1 P =  0 1 Now compute

 −1

P AP =  =

1 0 −1 0

−1 1 0 3

 

−1 0

4 3



1 0

1 1





which has the eigenvalues down the main diagonal, corresponding to the order in which the eigenvectors were written as columns of P. If we use the other order in writing the eigenvectors as columns, and define   1 1 Q=  1 0 then we get

 −1

Q AQ =

3 0

0 −1

 

9.2 Diagonalization of Matrices

333

Any linearly independent   eigenvectors   can be used in this diagonalization process. For 6 −4 example, if we use and , which are simply nonzero scalar multiples of the 0 −4 previously used eigenvectors, then we can define   6 −4 S= 0 −4 and now

 −1

S AS =

−1 0

0 3

 

EXAMPLE 9.7

Here is an example with more complicated arithmetic, but the idea remains the same. Let ⎞ ⎛ −1 1 3 4 ⎠ A=⎝ 2 1 1 0 −2 √ √ The eigenvalues are −1 − 21 + 21 29 − 21 − 21 29, with corresponding eigenvectors, respectively, √ √ ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 1 3+ √ 29 29 3− √ ⎝ −3 ⎠  ⎝ 10 + 2 29 ⎠  ⎝ 10 − 2 29 ⎠  1 2 2 These are linearly independent. Form the matrix √ ⎛ 1 3+ √ 29 P = ⎝ −3 10 + 2 29 1 2 We find that √ P−1 =

Then

⎛

232 √ 29

29 ⎜ √ ⎝ −16 − 2 29 √ 812 −16 − 2 29 ⎛ ⎜ P−1 AP = ⎜ ⎝

−1

√ ⎞ 3− √ 29 10 − 2 29 ⎠  2

− √116 29 √ −1 +√ 29 1 + 29

0

0

√ −1+ 29 2

0

0

⎞

√ ⎟ −19 + 5√ 29 ⎠  19 + 5 29

0 0

232 √ 29

√ −1− 29 2

⎞ ⎟ ⎟ ⎠

In this example, although we found P−1 explicitly, we did not actually need it to diagonalize A. Theorem 9.4 assures us that P−1 AP is a diagonal matrix with the eigenvalues of A down its main diagonal. All we really needed was to know that A had three linearly independent eigenvectors. This is a useful fact to keep in mind, particularly if P and P−1 are cumbersome to compute.

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EXAMPLE 9.8

Let



 −4 A=  −2 √ √ The eigenvalues are
−3 + 47i/2 and
−3 − 47i/2. Corresponding eigenvectors are, respectively,     8 8 √ √   1 − 47i 1 + 47i −1 3

Since these eigenvalues are linearly independent, there is a nonsingular 2 × 2 matrix P that diagonalizes A:  −3+√47i  0 2 −1 P AP =  √ −3− 47i 0 2 Of course, if we need P for some other calculation, as will occur later, we can write it down:   8 8 √ √ P=  1 − 47i 1 + 47i And, if we wish, we can compute P

−1

√

=

47i 752



√  8 −1 − i√47  1 − i 47 −8

However, even without explicitly writing P−1 , we know what P−1 AP is.

EXAMPLE 9.9

It is not necessary that A have n distinct eigenvalues in order to have n linearly independent eigenvectors. For example, let ⎛ ⎞ 5 −4 4 A = ⎝ 12 −11 12 ⎠  4 −4 5 The eigenvalues are 1 1 −3, with 1 having multiplicity 2. Associated with −3 we find an eigenvector ⎛ ⎞ 1 ⎝ 3 ⎠ 1 To find eigenvectors associated with 1 we must solve the system ⎞ ⎛ ⎞ ⎛ ⎞⎛ 0 −4 4 −4 x1
I3 − AX = ⎝ −12 12 −12 ⎠ ⎝ x2 ⎠ = ⎝ 0 ⎠  0 −4 4 −4 x3 This system has general solution

⎞ ⎞ ⎛ 0 1 ⎝ 0 ⎠+⎝ 1 ⎠ 1 −1 ⎛

9.2 Diagonalization of Matrices

335

We can therefore find two linearly independent eigenvectors associated with eigenvalue 1, for example, ⎛ ⎞ ⎛ ⎞ 1 0 ⎝ 0 ⎠ and ⎝ 1 ⎠  −1 1 We can now form the nonsingular matrix ⎛

1 1 0 P=⎝ 3 1 −1

⎞ 0 1 ⎠ 1

that diagonalizes A: ⎛

−3 P−1 AP = ⎝ 0 0

0 1 0

⎞ 0 0 ⎠ 1

Here is a proof of Theorem 9.4, explaining why a matrix P formed from linearly independent eigenvectors must diagonalize A. The proof makes use of an observation we have made before. When multiplying two n × n matrices A and B, column j of AB = A(column j of B). Let the eigenvalues of A be 1      n and corresponding eigenvectors V1      Vn . These form the columns of P. Since these eigenvectors are assumed to be linearly independent, the dimension of the column space of P is n. Therefore rank
P = n and P is nonsingular by Theorems 7.15 and 7.26. Now compute P−1 AP as follows. First,

Proof

column j of AP = A(column j of P = AVj = j Vj  Thus the columns of AP are 1 V1      n Vn and AP has the form ⎛ ⎞   ···  AP = ⎝ 1 V1 2 V2 · · · n Vn ⎠    ···  Then column j of P−1 AP = P−1 (column j of AP) = P−1 j Vj  = j P−1 Vj  But Vj is column j of P, so ⎛ ⎜ ⎜ ⎜ ⎜ P−1 Vj = column j of P−1 P = ⎜ ⎜ ⎜ ⎜ ⎝

0 0   1   0

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

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in which the column matrix on the right has all zero elements except 1 in row j. Combining the last two equations, we have ⎞ ⎛ ⎛ ⎞ 0 0 ⎜ 0 ⎟ ⎜ 0 ⎟ ⎟ ⎜ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟ ⎟ column j of P−1 AP = j P−1 Vj = j ⎜ ⎜ 1 ⎟ = ⎜ j ⎟  ⎟ ⎜ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟   ⎝  ⎠ ⎝  ⎠ 0 0 We now know the columns of P−1 AP, and ⎛ 1 ⎜ 0 ⎜ ⎜ P−1 AP = ⎜ 0 ⎜  ⎝  0

putting them together gives us ⎞ 0 0 ··· 0 2 0 · · · 0 ⎟ ⎟ 0 3 · · · 0 ⎟ ⎟     ⎟     ⎠ 0 0 · · · n

We can strengthen the conclusions of Theorem 9.4. So far, if A has n linearly independent eigenvectors, then we can diagonalize A. We will now show that this is the only time A can be diagonalized. Further, if, for any Q, Q−1 AQ is a diagonal matrix, then Q must have linearly independent eigenvectors of A as its columns.

THEOREM 9.5

Let A be an n × n diagonalizable matrix. Then A has n linearly independent eigenvectors. Further, if Q−1 AQ is a diagonal matrix, then the diagonal elements of Q−1 AQ are the eigenvalues of A, and the columns of Q are corresponding eigenvectors. Proof

Suppose that ⎛ ⎜ ⎜ Q−1 AQ = ⎜ ⎝

d1

O d2 

O



⎞ ⎟ ⎟ ⎟ = D ⎠

dn

Denote column j of Q as Vj . Then V1      Vn are linearly independent, because Q is nonsingular. We will show that dj is an eigenvalue of A, with corresponding eigenvector Vj . Write AQ = QD and compute both sides of this product separately. First, since the columns of Q are V1      Vn , ⎛ ⎞   ···  QD = ⎝ V1 V2 · · · Vn ⎠ D   ···  ⎛ ⎞   ···  = ⎝ d1 V1 d2 V2 · · · dn Vn ⎠    ··· 

9.2 Diagonalization of Matrices

337

a matrix having dj Vj as column j. Now compute ⎛ ⎞   ···  AQ = A ⎝ V1 V2 · · · Vn ⎠   ···  ⎛ ⎞   ···  = ⎝ AV1 AV2 · · · AVn ⎠    ···  a matrix having AVj as column j. Since AQ = QD, then column j of AQ equals column j of QD, so AVj = dj Vj which proves that dj is an eigenvalue of A with associated eigenvector Vj . As a consequence of this theorem, we see that not every matrix is diagonalizable.

EXAMPLE 9.10

Let

 B=

1 0

−1 1

 

B has eigenvalues 1 1, and every eigenvector has the form 



1 0

 . There are not two linearly

independent eigenvectors, so B is not diagonalizable. We could also proceed here by contradiction. If B were diagonalizable, then for some P,   1 0 −1 P AP =  0 1 From Theorem 9.5, the columns of P must be eigenvectors, so P must have the form     P=  0 0 But this matrix is singular (it has zero determinant, and its columns are multiples of each other, hence linearly dependent). Thus no matrix can diagonalize B. The key to diagonalization of an n × n matrix A is therefore the existence of n linearly independent eigenvectors. We saw (Example 9.9) that this does not require that the eigenvalues be distinct. However, if A does have n distinct eigenvalues, we claim that it must have n linearly independent eigenvectors, hence must be diagonalizable. THEOREM 9.6

Let the n×n matrix A have n distinct eigenvalues. Then corresponding eigenvectors are linearly independent. We will show by induction that any k distinct eigenvalues have associated with them k linearly independent eigenvectors. For k = 1, an eigenvector associated with a single eigenvalue is linearly independent, being a nonzero vector. Now suppose that any k − 1 distinct eigenvalues have associated with them k − 1 linearly independent eigenvectors. Suppose we have distinct

Proof

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eigenvalues 1      k . Let V1      Vk be associated eigenvectors. We want to show that V1      Vk are linearly independent. If these eigenvectors were linearly dependent, there would be numbers c1      ck not all zero such that c1 V1 + · · · + ck Vk = O By relabeling if necessary, we may suppose for convenience that c1 = 0. Now
1 In − A
c1 V1 + · · · + ck Vk  = O = c1
1 In − AV1 + c2
1 In − AV2 + · · · + ck
1 In − AVk = c1
1 V1 − AV1  + c2
1 V2 − AV2  + · · · + ck
1 Vk − AVk  = c1
1 V1 − 1 V1  + c2
1 V2 − 2 V2  + · · · + ck
1 Vk − k Vk  = c2
1 − 2 V2 + · · · + ck
1 − k Vk  But V2      Vk are linearly independent by the inductive hypothesis, so each of these coefficients must be zero. Since 1 − j = 0 for j = 2    k by the assumption that the eigenvalues are distinct, then c2 = · · · = ck = 0 But then c1 V1 = O. Since V1 is an eigenvector and cannot be O, then c1 = 0 also, a contradiction. Therefore V1      Vk are linearly independent.

COROLLARY 9.1

If an n × n matrix A has n distinct eigenvalues, then A is diagonalizable.

EXAMPLE 9.11

Let ⎛

−2 ⎜ 1 A=⎜ ⎝ 0 2 The eigenvalues of A are 3 4 − 25 + diagonalizable. For some P, ⎛

3 ⎜ 0 ⎜ P−1 AP = ⎜ ⎝ 0 0

1 2

0

0 0 4 0

⎞ 5 0 ⎟ ⎟ 0 ⎠ −3

√ √ 41 − 25 − 21 41. Because these are distinct, A is

0 4 0

0 3 4 0

0 0

− 25

0 0

√ + 41 1 2

0

0

− 25

√ − 21 41

We do not need to actually produce P explicitly to conclude this.

⎞ ⎟ ⎟ ⎟ ⎠

9.3 Orthogonal and Symmetric Matrices

SECTION 9.2

PROBLEMS ⎛

−2 ⎜ −4 10. ⎜ ⎝ 0 0

In each of Problems 1 through 10, produce a matrix that diagonalizes the given matrix, or show that this matrix is not diagonalizable.   0 −1 1. 4 3   5 3 2. 1 3   1 0 3. −4 1   −5 3 4. 0 9 ⎞ ⎛ 5 0 0 3 ⎠ 5. ⎝ 1 0 0 0 −2 ⎞ ⎛ 0 0 0 6. ⎝ 1 0 2 ⎠ 0 1 3 ⎞ ⎛ −2 0 1 0 ⎠ 7. ⎝ 1 1 0 0 −2 ⎞ ⎛ 2 0 0 2 1 ⎠ 8. ⎝ 0 0 −1 2 ⎞ ⎛ 1 0 0 0 ⎜ 0 4 1 0 ⎟ ⎟ 9. ⎜ ⎝ 0 0 −3 1 ⎠ 0 0 1 −2

9.3

339

0 −2 0 0

⎞ 0 0 ⎟ ⎟ 0 ⎠ −2

0 0 −2 0

11. Suppose A2 is diagonalizable. Prove that A is diagonalizable. 12. Let A have eigenvalues 1      n and suppose P diagonalizes A. Prove that, for any positive integer k, ⎛ ⎜ ⎜ Ak = P ⎜ ⎝

k1

O k2 

O



⎞ ⎟ ⎟ −1 ⎟P  ⎠

kn

In each of Problems 13 through 16, compute the indicated power of the matrix, using the idea of Problem 12.   −1 0  A18 13. A = 1 −5   −3 −3  A16 14. A = −2 4   0 −2  A43 15. A = 1 0   −2 3  A31 16. A = 3 −4

Orthogonal and Symmetric Matrices Recall that the transpose of a matrix is obtained by interchanging the rows with the columns. For example, if  A=

−6 1

3 −7

−6 3

1 −7



then  At =

 

Usually At is simply another matrix. However, in the special circumstance that the transpose of a matrix is its inverse, we call A an orthogonal matrix.

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DEFINITION 9.5

Orthogonal Matrix

A square matrix A is orthogonal if and only if AAt = At A = In .

An orthogonal matrix is therefore nonsingular, and we find its inverse simply by taking its transpose.

EXAMPLE 9.12

Let

⎛

0

⎜ A=⎝ 1 0 Then

⎛

0

⎜ AAt = ⎝ 1 0

√1 5

√2 5

0

0

√2 5

− √15

√1 5

√2 5

0

0

√2 5

− √15

⎞⎛ ⎟⎜ ⎠⎝

⎞ ⎟ ⎠

0

1

0

√1 5 √2 5

0

√2 5 − √15

0

⎞ ⎟ ⎠ = I3

and a similar calculation gives At A = I3 . Therefore this matrix is orthogonal, and ⎞ ⎛ 0 1 0 ⎟ ⎜ √2 A−1 = At = ⎝ √15 0 ⎠ 5 √2 √1 0 − 5 5 Because the transpose of the transpose of a matrix is the original matrix, a matrix is orthogonal exactly when its transpose is orthogonal. THEOREM 9.7

A is an orthogonal matrix if and only if At is an orthogonal matrix. Orthogonal matrices have several interesting properties. We will show first that the determinant of an orthogonal matrix must be 1 or −1. THEOREM 9.8

If A is an orthogonal matrix, then A = ±1.  t 2 t Proof Since AA = In , then AA  = 1 = A At  = A . The next property of orthogonal matrices is actually the rationale for the name orthogonal. A set of vectors in Rn is said to be orthogonal if any two distinct vectors in the set are orthogonal (that is, their dot product is zero). The set is orthonormal if, in addition, each vector has length 1. We claim that the rows of an orthogonal matrix form an orthonormal set of vectors, as do the columns.

9.3 Orthogonal and Symmetric Matrices

341

This can be seen in the matrix of the last example. The row vectors are     0 √15 √25  1 0 0  0 √25 − √15  These each have length 1, and each is orthogonal to each of the other two. Similarly, the columns of that matrix are ⎛ ⎞ ⎛ √1 ⎞ ⎛ √2 ⎞ 5 5 0 ⎟ ⎜ ⎟ ⎝ 1 ⎠⎜ ⎝ 0 ⎠⎝ 0 ⎠ 0 √2 − √1 5

5

Each is orthogonal to the other two, and each has length 1. Not only do the row (column) vectors of an orthogonal matrix form an orthonormal set of vectors in Rn , but this property completely characterizes orthogonal matrices. THEOREM 9.9

Let A be a real n × n matrix. Then 1. A is orthogonal if and only if the row vectors form an orthonormal set of vectors in Rn . 2. A is orthogonal if and only if the column vectors form an orthonormal set of vectors in Rn . Proof Recall that the i j element of AB is the dot product of row i of A with column j of B. Further, the columns of At are the rows of A. Therefore,

i j element of AAt = (row i of A ·
column j of At ) =
row i of A) ·
row j of A). Now suppose that A is an orthogonal matrix. Then AAt = In , so the i j element of AAt is zero if i = j. Therefore the dot product of two distinct rows of A is zero, and the rows form an orthogonal set of vectors. Further, the dot product of row i with itself is the i i element of AAt , and this is 1, so the rows form an orthonormal set of vectors. Conversely, suppose the rows of A form an orthonormal set of vectors. Then the dot product row i with row j is zero if i = j, so the i j element of AAt is zero if i = j. Further, the i i element of AAt is the dot product of row i with itself, and this is 1. Therefore AAt = In . Similarly, At A is In , so A is an orthogonal matrix. This proves (1). A proof of (2) is similar. We now have a great deal of information about orthogonal matrices. We will use this to completely determine all 2 × 2 orthogonal matrices. Let   a b Q=  c d What do we have to say about a, b, c and d to make this an orthogonal matrix? First, the two row vectors must be orthogonal (zero dot product), and must have length 1, so ac + bd = 0

(9.1)

a2 + b 2 = 1

(9.2)

c2 + d2 = 1

(9.3)

The two column vectors must also be orthogonal, so in addition, ab + cd = 0

(9.4)

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Finally, Q = ±1, so ad − bc = ±1 This leads to two cases. Case 1—ad − bc = 1. Multiply equation (9.1) by d to get acd + bd2 = 0 Substitute ad = 1 + bc into this equation to get c
1 + bc + bd2 = 0 or c + b
c2 + d2  = 0 But c2 + d2 = 1 from equation (9.3), so c + b = 0, hence c = −b Put this into equation (9.4) to get ab − bd = 0 Then b = 0 or a = d, leading to two subcases. Case 1-(a)—b = 0. Then c = −b = 0 also, so   a 0 Q=  0 d But each row vector has length 1, so a2 = d2 = 1. Further, Q = ad = 1 in the present case, so a = d = 1 or a = d = −1. In these cases, Q = I2 Case 1-(b)—b = 0. Then a = d, so Q=

or 

a −b

Q = −I2 

b a

 

Since a2 + b2 = 1, there is some  in 0 2 such that a = cos
 and b = sin
. Then   cos
 sin
 Q=  − sin
 cos
 This includes the two results of case 1(a) by choosing  = 0 or  = . Case 2—ad − bc = −1. By an analysis similar to that just done, we find now that, for some ,   cos
 sin
 Q=  sin
 − cos
 These two cases give all the 2 × 2 orthogonal matrices. For example, with  = /4 we get the orthogonal matrices    1  1 1 1 √ 2 − √12

√

2 √1 2

and

√ 2 √1 2

√ 2 − √12

9.3 Orthogonal and Symmetric Matrices and with  = /6 we get



√ 3 2

− 21

1 2 √ 3 2



 and

√ 3 2 1 2

1 2 √ − 23

343

 

We can recognize the orthogonal matrices   cos
 sin
 − sin
 cos
 as rotations in the plane. If the positive x y system is rotated counterclockwise  radians to form a new x
 y
system, the coordinates in the two systems are related by 
    x cos
 sin
 x =  y
− sin
 cos
 y We will now consider another kind of matrix that is related to the class of orthogonal matrices.

DEFINITION 9.6

Symmetric Matrix

A square matrix is symmetric if A = At .

This means that each aij = aji , or that the main diagonal. For example, ⎛ −7 ⎜ −2 ⎜ ⎝ 1 14

the matrix elements are the same if reflected across −2 2 −9 47

⎞ 1 14 −9 47 ⎟ ⎟ 6  ⎠  22

is symmetric. A symmetric matrix need not have real numbers as elements. However, when it does, it has the remarkable property of having only real eigenvalues. THEOREM 9.10

The eigenvalues of a real, symmetric matrix are real numbers. Before showing why this is true, we will review √ some facts about complex numbers. A complex number z = a + ib has magnitude z = a2 + b2 . The conjugate of z is defined to be z = a − ib. When z is represented as the point
a b in the plane, z is the point
a −b, which is the reflection of
a b across the x-axis. A number is real exactly when it equals its own conjugate. Further, zz = a2 + b2 = z2 and
z = z We take the conjugate A of a matrix A by taking the conjugate of each of its elements. The product of a conjugate is the conjugate of a product:   
AB = A B 

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Further, the operation of taking the conjugate commutes with the operation of taking a transpose: t

C =
Ct  For example, if ⎛

⎞ 1 − 2i 0 ⎠ 4

⎛

⎞ 1 + 2i 0 ⎠ 4

i 3 C=⎝ −2 + i then

−i 3 C=⎝ −2 − i and t

C =



−i 1 + 2i

−2 − i 4

3 0

 =
Ct 

We will now prove that the eigenvalues of a real symmetric matrix must be real. Proof

Let A be an n × n matrix of real numbers. Let  be an eigenvalue, and let ⎛ ⎞ e1 ⎜ e2 ⎟ ⎜ ⎟ E=⎜  ⎟ ⎝  ⎠ en

be an associated eigenvector. Then AE = E. Multiply this equation on the left by the 1 × n matrix   t E = e1 e2 · · · en to get t

t

t

E AE = E E = E E

=



e1

e2

···

⎛ en

⎜ ⎜ ⎜ ⎝

e1 e2  

⎞ ⎟ ⎟ ⎟ ⎠

en =  e1 e1 + e2 e2 + · · · + en en   =  e1 2 + e2 2 + · · · + en 2 

(9.5)

which is a real number. Here we are using the standard convention that a 1 × 1 matrix is identified with its single element. Now compute  t t E AE = E A E = Et AE in which we have used the fact that A has real elements to write A = A.

(9.6)

9.3 Orthogonal and Symmetric Matrices

345

Now Et AE is a 1×1 matrix, and so is the same as its transpose. Recalling that the transpose of a product is the product of the transposes in the reverse order, take the transpose of the last equation (9.6) to get  t t t Et AE = Et AE = E A
Et t = E AE (9.7) From equations (9.6) and (9.7) we have t

t

E AE = E AE t

Therefore the 1 × 1 matrix E AE, being equal to its conjugate, is a real number. Now return to equation (9.5). We have just shown that the left side of this equation is real. Therefore the right side must be real. But e1 2 + e2 2 + · · · + en 2 is certainly real. Therefore  is real, and the theorem is proved. One ramification of this theorem is that a real, symmetric matrix also has real eigenvectors. We claim that, more than this, eigenvectors from distinct eigenvalues are orthogonal. THEOREM 9.11

Let A be a real symmetric matrix. Then eigenvectors associated with distinct eigenvalues are orthogonal. Proof

Let  and  be distinct eigenvalues with, respectively, eigenvectors ⎛ ⎞ ⎞ ⎛ g1 e1 ⎜ g2 ⎟ ⎜ e2 ⎟ ⎜ ⎟ ⎟ ⎜ E = ⎜  ⎟ and G = ⎜  ⎟  ⎝  ⎠ ⎝  ⎠ en gn

Identifying, as usual, a real number with the 1×1 matrix having this number as its only element, the dot product of these two n-vectors can be written as a matrix product e1 g1 + · · · + en gn = Et G Since AE = E and AG = G, we have Et G =
Et G =
AEt G =
Et At G =
Et AG = Et
AG = Et
G = Et G Then
 − Et G = 0 But  = , so Et G = 0 and the dot product of these two eigenvectors is zero. These eigenvectors are therefore orthogonal.

EXAMPLE 9.13

Let

⎛

3 0 A=⎝ 0 2 −2 0

⎞ −2 0 ⎠ 0

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a 3 × 3 real symmetric matrix. The eigenvalues are 2 −1 4, with associated eigenvectors ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 0 1 2 ⎝ 1 ⎠⎝ 0 ⎠⎝ 0 ⎠ 0 2 −1 These form an orthogonal set of vectors. In this example, the eigenvectors, while orthogonal to each other, are not all of length 1. However, a scalar multiple of an eigenvector is an eigenvector, so we can also write the following eigenvectors of A: ⎛ ⎞ ⎛ √1 ⎞ ⎛ √2 ⎞ 0 5 5 ⎟ ⎜ ⎟ ⎝ 1 ⎠⎜ ⎝ 0 ⎠⎝ 0 ⎠ 0 √2 − √15 5 These are still mutually orthogonal (multiplying by a positive scalar does not change orientation), but are now orthonormal. They can therefore be used as columns of an orthogonal matrix ⎛ ⎞ √2 0 √15 5 ⎜ ⎟ Q=⎝ 1 0 0 ⎠ 0 √25 − √15 These column vectors, being orthogonal to each other, are linearly independent by Theorem 6.14. But whenever we form a matrix from linearly independent eigenvectors of A, this matrix diagonalizes A. Further, since Q is an orthogonal matrix, Q−1 = Qt . Therefore, as we can easily verify in this example, ⎛ ⎞ 2 0 0 Q−1 AQ = ⎝ 0 −1 0 ⎠  0 0 4 The idea we have just illustrated forms the basis for the following result. THEOREM 9.12

Let A be a real, symmetric matrix. Then there is a real, orthogonal matrix that diagonalizes A.

EXAMPLE 9.14

Let

⎛

2 A=⎝ 1 0 The eigenvalues are

√

1 −2 4

⎞ 0 4 ⎠ 2

√ 21 − 21 and 2, with associated eigenvectors, respectively ⎞ ⎞ ⎛ ⎞ ⎛ ⎛ −4 √ 1 √1 ⎝ 21 − 2 ⎠  ⎝ − 21 − 2 ⎠  ⎝ 0 ⎠  1 4 4

9.4 Quadratic Forms

347

These eigenvectors are mutually orthogonal, but not orthonormal. Divide each eigenvector by its length to get the three new eigenvectors: ⎞ ⎞ ⎛ ⎛ ⎛ ⎞ 1 −4 √1 1⎝ √ 1 1 21 − 2 ⎠  ⎝ − 21 − 2 ⎠  ⎝ 0 ⎠    1 4 4 where

" =

√ 42 − 4 21

and  =

√

17

The orthogonal matrix Q having these normalized eigenvectors as columns diagonalizes A.

SECTION 9.3

PROBLEMS

In each of Problems 1 through 12, find the eigenvalues of the matrix and, for each eigenvalue, a corresponding eigenvector. Check that eigenvectors associated with distinct eigenvalues are orthogonal. Find an orthogonal matrix that diagonalizes the matrix.   4 −2 1. −2 1   −3 5 2. 5 4   6 1 3. 1 4   −13 1 4. 1 4 ⎛ ⎞ 0 1 0 5. ⎝ 1 −2 0 ⎠ 0 0 3 ⎛ ⎞ 0 1 1 6. ⎝ 1 2 0 ⎠ 1 0 2

9.4

⎛

5 7. ⎝ 0 2 ⎛ 8.

9.

10.

11.

12.

2 ⎝ −4 0 ⎛ 0 ⎝ 0 0 ⎛ 1 ⎝ 3 0 ⎛ 0 ⎜ 0 ⎜ ⎝ 0 0 ⎛ 5 ⎜ 0 ⎜ ⎝ 0 0

⎞ 2 0 ⎠ 0

0 0 0

−4 0 0 0 1 −2 3 0 1

⎞ 0 0 ⎠ 0 ⎞ 0 −2 ⎠ 0 ⎞

0 1 ⎠ 1

0 1 −2 0

0 −2 1 0

0 0 −1 0

0 −1 0 0

⎞ 0 0 ⎟ ⎟ 0 ⎠ 0 ⎞ 0 0 ⎟ ⎟ 0 ⎠ 0

Quadratic Forms DEFINITION 9.7

A (complex) quadratic form is an expression n n  

ajk zj zk 

j=1 k=1

in which each ajk and zj are complex numbers.

(9.8)

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For n = 2 this quadratic form is a11 z1 z1 + a12 z1 z2 + a21 z1 z2 + a22 z2 z2  The terms involving zj zk with j = k, are the mixed product terms. The quadratic form is real if each ajk and zj is real. In this case we usually write zj as xj . Since xj = xj when xj is real, the form (9.8) in this case is n n  

ajk xj xk 

j=1 k=1

For n = 2, this is a11 x12 +
a12 + a21 x1 x2 + a22 x22  The terms involving x12 and x22 are the squared terms in this real quadratic form, and x12 is the mixed product term. It is often convenient to write a quadratic form (9.8) in matrix form. If A = aij  and ⎛ ⎜ ⎜ Z=⎜ ⎝

⎞

z1 z2  

⎟ ⎟ ⎟ ⎠

zn then ⎛ t

Z AZ =

=





z1

z2

···

zn

⎜ ⎜ ⎜ ⎝

a11 z1 + · · · + a1n zn

a11 a21  

a12 a22  

··· ···  

a1n a2n  

an1

an2

···

ann

···

⎞⎛ ⎟⎜ ⎟⎜ ⎟⎜ ⎠⎝

z1 z2  

⎞ ⎟ ⎟ ⎟ ⎠

zn ⎛ z1 z ⎜ ⎜ 2 an1 z1 + · · · + ann zn ⎜  ⎝ 

⎞ ⎟ ⎟ ⎟ ⎠

zn = a11 z1 z1 + · · · + a1n zn z1 + · · · + an1 z1 zn + · · · + ann zn zn =

n  n 

ajk zj zk 

j=1 k=1

Similarly, any real quadratic form can be written in matrix form as Xt AX. Given a quadratic form, we may choose different matrices A such that the form is Zt AZ.

EXAMPLE 9.15

Let  A=

1 3

4 2

 

9.4 Quadratic Forms Then 

x1

x2





1 3

4 2





x1 x2

=



x1 + 3x2

4x1 + 2x2





x1 x2

349



= x12 + 3x1 x2 + 4x1 x2 + 2x22 = x12 + 7x1 x2 + 2x22  But we can also write this quadratic form as x12 +

 7 7 x1 x2 + x2 x1 + 2x22 = x1 2 2

x2





1

7 2

7 2

2



x1 x2

 

The advantage of the latter formulation is that the quadratic form is Xt AX with A a symmetric matrix. There is an expression involving a quadratic form that gives the eigenvalues of a matrix in terms of an associated eigenvector. We will have use for this shortly.

LEMMA 9.1

Let A be an n × n matrix of real or complex numbers. Let  be an eigenvalue with eigenvector Z. Then t

=

Proof

t

Z AZ t



ZZ

t

Since AZ = Z, then Z AZ = Z Z.

Using a calculation done in equation (9.5), we can write n  n  1  =   2 ajk zj zk  n   j=1 k=1 j=1 zj

Quadratic forms arise in a variety of contexts. In mechanics, the kinetic energy of a particle is a real quadratic form, and in analytic geometry a conic is the locus of points in the plane for which a quadratic form in the coordinates is equal to some constant. For example, 1 x12 + x22 = 9 4 is the equation of an ellipse in the x1  x2 plane. In some problems involving quadratic forms, calculations are simplified if we transform from the x1  x2      xn coordinate system to a y1  y2      yn system in which there are no mixed product terms. That is, we want to choose y1      yn so that n n   j=1 k=1

aij xj xk =

n 

j yj2 

j=1

The y1      yn coordinates are called principal axes for the quadratic form.

(9.9)

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This kind of transformation is commonly done in analytic geometry, where a rotation of axes is used to eliminate mixed product terms in the equation of a conic. For example, the change of variables 1 1 x1 = √ y1 + √ y2 2 2 1 1 x2 = √ y1 − √ y2 2 2 transforms the quadratic form x12 − 2x1 x2 + x22 to 2y22  with no mixed product term. Using this transformed form, we could analyze the graph of x12 − 2x1 x2 + x22 = 4 in the x1  x2 − system, in terms of the graph of y22 = 2 in the y1  y2 − system. √ In the y1  y2 − plane it is clear that the graph consists of two horizontal straight lines y2 = ± 2. We will now show that a transformation that eliminates the mixed product terms of a real quadratic form always exists. THEOREM 9.13

Principal Axis Theorem

Let A be a real symmetric matrix with eigenvalues 1      n . Let Q be orthogonal matrix  an  that diagonalizes A. Then the change of variables X = QY transforms nj=1 nk=1 ajk xj xk to n 

j yj2 

j=1

Proof

The proof is a straightforward calculation: n  n 

aij xj xk = Xt AX

j=1 k=1

=
QYt A
QY =
Yt Qt A
QY = Yt
Qt AQY  = y1

···

⎛

⎜ ⎜ yn ⎜ ⎝

O

1 2  O

 n

⎞

⎛ ⎟ ⎟⎜ ⎟⎝ ⎠

⎞ y1  ⎟  ⎠ yn

= 1 y12 + · · · + n yn2  The expression 1 y12 + · · · + n yn2 is called the standard form of the quadratic form Xt AX.

9.4 Quadratic Forms

351

EXAMPLE 9.16

Consider again x12 − 2x1 x2 + x22 . This is Xt AX with   1 −1 A=  −1 1 The eigenvalues of A are 0 and 2, with corresponding eigenvectors     1 1 and  1 −1 Dividing each eigenvector by its length, we obtain the eigenvectors  1   1  √

2 √1 2

√ 2 − √12

and



These form the columns of an orthogonal matrix Q that diagonalizes A:   1 1 Q= The transformation defined by X = QY is    x1 = x2

√ 2 √1 2

√1 2 √1 2

√

2

− √12

√1 2 − √12





y1 y2

 

which gives exactly the transformation used above to reduce the quadratic form x12 − 2x1 x2 + x22 to the standard form 2y22 .

EXAMPLE 9.17

Analyze the conic 4x12 − 3x1 x2 + 2x22 = 8. First write the quadratic form as Xt AX = 8, where   4 − 23 A=  − 23 2 √ The eigenvalues of A are
6 ± 13/2. By the principal axis theorem there is an orthogonal matrix Q that transforms the equation of the conic to standard form: √ √ 6 + 13 2 6 − 13 2 y1 + y2 = 8 2 2 This is an ellipse in the y1  y2 − plane. Figure 9.2 shows a graph of this ellipse. x2 y 2

x1 y1 FIGURE 9.2

CHAPTER 9

352

Eigenvalues, Diagonalization, and Special Matrices

PROBLEMS

SECTION 9.4

In each of Problems 1 through 6, find a matrix A such that the quadratic form is Xt AX. 1.

x12 + 2x1 x2 + 6x22

14. x12 − 2x1 x2 + 4x22 = 6 15. 3x12 + 5x1 x2 − 3x22 = 5 16. −2x12 + 3x22 + x1 x2 = 5

2. 3x12 + 3x22 − 4x1 x2 − 3x1 x3 + 2x2 x3 + x32

17. 4x12 − 4x22 + 6x1 x2 = 8

3. x12 − 4x1 x2 + x22

18. 6x12 + 2x1 x2 + 5x22 = 14

4.

2x12 − x22 + 2x1 x2

5. −x12 + x42 − 2x1 x4 + 3x2 x4 − x1 x3 + 4x2 x3 6. x12 − x22 − x1 x3 + 4x2 x3 In Problems 7 through 13, find the standard form of the quadratic form. 7. −5x12 + 4x1 x2 + 3x22 8. 4x12 − 12x1 x2 + x22 9. −3x12 + 4x1 x2 + 7x22 10. 4x12 − 4x1 x2 + x22 11. −6x1 x2 + 4x22 12. 5x12 + 4x1 x2 + 2x22 13. −2x1 x2 + 2x32 In each of Problems 14 through 18, use the principal axis theorem to analyze the conic.

9.5

In each of Problems 19 through 22, write the quadratic form defined by the matrix.   −2 1 19. 1 6 ⎞ ⎛ 14 −3 0 2 1 ⎠ 20. ⎝ −3 0 1 7 ⎞ ⎛ 6 1 −7 0 ⎠ 21. ⎝ 1 2 −7 0 1 ⎞ ⎛ 7 1 −2 0 −1 ⎠ 22. ⎝ 1 −2 −1 3 23. Give an example of a real, 3 × 3 matrix that cannot be the coefficient matrix of a real quadratic form.

Unitary, Hermitian, and Skew Hermitian Matrices If U is a nonsingular complex matrix, then U−1 exists and is generally also a complex matrix. We claim that the operations of taking the complex conjugate and of taking a matrix inverse can be performed in either order. LEMMA 9.2

U

−1

= U−1 .

Proof

We know that the conjugate of a product is the product of the conjugates, so In = In = UU−1 = U U−1 

This implies that U−1 is the inverse of U. Now define a matrix to be unitary if the inverse of its conjugate (or conjugate of its inverse) is equal to its transpose.

DEFINITION 9.8

Unitary Matrix

An n × n complex matrix U is unitary if and only if U

−1

= Ut .

9.5 Unitary, Hermitian, and Skew Hermitian Matrices

353

This condition is equivalent to saying that UUt = In .

EXAMPLE 9.18

Let

 U=

√ i/√2 −i/ 2

√  1/√2  1/ 2

Then U is unitary because

√ √  √ −i/√2 1/√2 i/√2 UU = i/ 2 1/ 2 1/ 2   1 0 =  0 1

√  −i/√2 1/ 2



t

If U is a real matrix, then the unitary condition UUt = In becomes UUt = In , which makes U an orthogonal matrix. Unitary matrices are the complex analogues of orthogonal matrices. Since the rows (or columns) of an orthogonal matrix form an orthonormal set of vectors, we will develop the complex analogue of the concept of orthonormality. Recall that, for two vectors
x1      xn  and
y1     yn  in Rn , we can define the column matrices ⎛ ⎛ ⎞ ⎞ y1 x1 ⎜ y2 ⎟ ⎜ x2 ⎟ ⎜ ⎜ ⎟ ⎟ X = ⎜  ⎟ and Y = ⎜  ⎟   ⎝  ⎠ ⎝  ⎠ xn

yn

and obtain the dot product X · Y as Xt Y. In particular, this gives the square of the length of X as Xt X = x12 + x22 + · · · + xn2  To generalize this to the complex case, suppose we have complex n-vectors
z1  z2      zn  and
w1  w2      wn . Form the column matrices ⎛ ⎛ ⎞ ⎞ w1 z1 ⎜ w2 ⎟ ⎜ z2 ⎟ ⎜ ⎜ ⎟ ⎟ Z = ⎜  ⎟ and W = ⎜  ⎟  ⎝  ⎠ ⎝  ⎠ zn

wn

It is tempting to define the dot product of these complex vectors as Zt W. The problem with this is that then we get Zt Z = z21 + z22 + · · · + z2n and this will in general be complex. We want to interpret the dot product of a vector with itself as the square of its length, and this should be a nonnegative real number. We get around this by defining the dot product of complex Z and W to be t

Z · W = Z W = z1 w1 + z2 w2 + · · · + zn wn 

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In this way the dot product of Z with itself is t

Z Z = z1 z1 + z2 z2 + · · · + zn zn = z1 2 + z2 2 + · · · + zn 2  a nonnegative real number With this as background, we will define the complex analogue of an orthonormal set of vectors.

DEFINITION 9.9

Unitary System of Vectors

Complex n-vectors F1      Fr form a unitary system if Fj · Fk = 0 for j = k, and each Fj · Fj = 1.

If each Fj has all real components, then this corresponds exactly to an orthonormal set of vectors in Rn . We can now state the analogue of Theorem 9.9 for unitary matrices. THEOREM 9.14

Let U be an n × n complex matrix. Then U is unitary if and only if its row vectors form a unitary system. The proof is like that of Theorem 8.9, and is left to the student. It is not difficult to show that U is also unitary if and only if its column vectors form a unitary system.

EXAMPLE 9.19

Consider again  U=

√ i/√2 −i/ 2

√  1/√2  1/ 2

The row vectors, written as 2 × 1 matrices, are √   √   i/√2 −i/√2 F1 = and F2 =  1/ 2 1/ 2 Then F1 · F2 = F1 · F1 =

 

√

i/ 2 √ i/ 2

√  −i/√2 = 0 1/ 2 1/ 2  √  √  i/ 2 √ =1 1/ 2 1/ 2 √



√



and F2 · F2 =



√

−i/ 2

1/ 2





√  −i/√ 2 = 1 1/ 2

9.5 Unitary, Hermitian, and Skew Hermitian Matrices

355

We will show that the eigenvalues of a unitary matrix must lie on the unit circle in the complex plane.

THEOREM 9.15

Let  be an eigenvalue of the unitary matrix U. Then  = 1. Let E be an eigenvector associated with . Then UE = E, so UE = E. Then  t t UE = E 

Proof

so t

t

t

E U = E  t

But U is unitary, so U = U−1 , and t

t

E U−1 = E  Multiply both sides of this equation on the right by UE to get t

t

t

t

E E =  E UE =  E E = E E t

Now E E is the dot product of the eigenvector with itself, and so is a positive number. Dividing t the last equation by E E gives  = 1. But then 2 = 1, so  = 1. We have defined a matrix to be unitary if its transpose is the conjugate of its inverse. A matrix is hermitian if its transpose is equal to its conjugate. If the transpose equals the negative of its conjugate, the matrix is called skew-hermitian.

DEFINITION 9.10

1. Hermitian Matrix An n × n complex matrix H is hermitian if and only if H = Ht . 2. Skew-Hermitian Matrix An n × n complex matrix S is skew-hermitian if and only if S = −St .

In the case that H has real elements, hermitian is the same as symmetric, because in this case H = H.

EXAMPLE 9.20

Let

⎛

15 H = ⎝ −8i 6 + 2i

8i 0 −4 − i

⎞ 6 − 2i −4 + i ⎠  −3

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Then ⎛

⎞ 6 + 2i −4 − i ⎠ = Ht  −3

−8i 0 −4 + i

15 H = ⎝ 8i 6 − 2i so H is hermitian. If ⎛

0 S = ⎝ 8i 2i then S is skew-hermitian because ⎛

−8i 0 −4i

0 S = ⎝ −8i −2i

8i 0 4i

⎞ 2i 4i ⎠ 0

⎞ −2i −4i ⎠ = −St  0

The following theorem says something about quadratic forms with hermitian or skewhermitian matrices.

THEOREM 9.16

Let ⎛ ⎜ ⎜ Z=⎜ ⎝

z1 z2  

⎞ ⎟ ⎟ ⎟ ⎠

zn be a complex matrix. Then t

1. If H is hermitian, then Z HZ is real. t 2. If S is skew-hermitian, then Z SZ is zero or pure-imaginary. t

For (1), suppose H is hermitian. Then H = H, so

Proof

t

t

Z HZ = Z HZ = Zt HZ t

But Z HZ is a 1 × 1 matrix and so equals its own transpose. Continuing from the last equation, we have  t t t t Zt HZ = Zt HZ = Z H
Zt t = Z HZ Therefore t

t

Z HZ = Z HZ t

t

Since Z HZ equals its own conjugate, then Z HZ is real.

9.5 Unitary, Hermitian, and Skew Hermitian Matrices

357

t

To prove (2) suppose S is skew-hermitian. Then S = −S. By an argument like that done in the proof of (1), we get t

t

Z SZ = −Z SZ t

Now write Z SZ =  + i. The last equation becomes  − i = − − i t

Then  = −, so  = 0 and Z SZ is pure imaginary. Using these results on quadratic forms, we can say something about the eigenvalues of hermitian and skew-hermitian matrices. THEOREM 9.17

1. The eigenvalues of a hermitian matrix are real. 2. The eigenvalues of a skew-hermitian matrix are zero or pure imaginary. For (1), let  be an eigenvalue of the hermitian matrix H , with associated eigenvector E. By Lemma 8.1,

Proof

t

=

E HE t



EE

But by (1) of the preceding theorem, the numerator of this quotient is real. The denominator is the square of the length of E, and so is also real. Therefore  is real. For (2), let  be an eigenvalue of the skew-hermitian matrix S, with associated eigenvector E. Again by Lemma 8.1, t

=

E SE t



EE

By (2) of the preceding theorem, the numerator of this quotient is either zero or pure imaginary. Since the denominator is a positive real number, then  is either zero or pure imaginary. Figure 9.3 shows a graphical representation of the conclusions of Theorems 9.15 and 9.17. When plotted as points in the complex plane, eigenvalues of a unitary matrix lie on the unit circle about the origin, eigenvalues of a hermitian matrix lie on the horizontal (real) axis, and eigenvalues of a skew-hermitian matrix lie on the vertical (imaginary) axis.

Imaginary axis Unitary (eigenvalues have magnitude 1)

Skew-hermitian (pure imaginary eigenvalues) Real axis

1

Hermitian (real eigenvalues) FIGURE 9.3 Eigenvalue locations.

358

CHAPTER 9

SECTION 9.5

Eigenvalues, Diagonalization, and Special Matrices

PROBLEMS

In each of Problems 1 through 9, determine whether the matrix is unitary, hermitian, skew-hermitian, or none of these. Find the eigenvalues of each matrix and an associated eigenvector for each eigenvalue. Determine which matrices are diagonalizable. If a matrix is diagonalizable, produce a matrix that diagonalizes it.   0 2i 1. 2i 4   3 4i 2. 4i −5 ⎞ ⎛ 0 1 0 0 1−i ⎠ 3. ⎝ −1 0 −1 − i 0 √ √ ⎛ ⎞ 1/ √2 i/√2 0 4. ⎝ −1/ 2 i/ 2 0 ⎠ 0 0 1 ⎞ ⎛ 3 2 0 0 i ⎠ 5. ⎝ 2 0 −i 0

⎛

−1 6. ⎝ 0 3+i ⎛ i 7. ⎝ −1 0 ⎛ 3i 8. ⎝ −1 0 ⎛ 8 9. ⎝ −1 −i

0 1 0 1 0 2i 0 0 −i −1 0 0

⎞ 3−i 0 ⎠ 0 ⎞ 0 2i ⎠ 0 ⎞ 0 i ⎠ 0 ⎞ i 0 ⎠ 0

10. Let A be unitary, hermitian or skew-hermitian. Prove that AAt = AA. 11. Prove that the main diagonal elements of a skewhermitian matrix must be zero or pure imaginary. 12. Prove that the main diagonal elements of a hermitian matrix must be real. 13. Prove that the product of two unitary matrices is unitary.

PA RT

3

CHAPTER 10 Systems of Linear Differential Equations

CHAPTER 11 Qualitative Methods and Systems of Nonlinear Differential Equations

Systems of Differential Equations and Qualitative Methods

We will now use matrices to study systems of differential equations. These arise, for example, in modeling mechanical and electrical systems having more than one component. We will separate our study of systems of differential equations into two chapters. The first, Chapter 10, is devoted to systems of linear differential equations. For these, powerful matrix methods can be brought to bear to write solutions. For systems of nonlinear differential equations, for which we usually cannot write explicit solutions, we must develop a different set of tools designed to determine qualitative properties of solutions. This is done in Chapter 11.

359

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CHAPTER

10

THEORY OF SYSTEMS OF LINEAR FIRST-ORDER DIFFERENTIAL EQUATIONS SOLUTION OF X
= AX WHEN A IS CONSTANT SOLUTION OF X
= AX + G THEORY OF SYSTEMS OF LINEAR FIRST-ORDER

Systems of Linear Differential Equations

Before beginning to study linear systems, recall from Section 2.6.4 that a linear differential equation of order n always gives rise to a system of n first-order linear differential equations, in such a way that the solution of the system gives the solution of the original nth order equation. Systems can be treated using matrix techniques, which are now at our disposal. For this reason we did not spend time on differential equations of order higher than 2 in Part 1. We will assume familiarity with vectors in Rn , matrix algebra, determinants, and eigenvalues and eigenvectors. These can be reviewed as needed from Part 2. We begin by laying the foundations for the use of matrices to solve linear systems of differential equations.

10.1

Theory of Systems of Linear First-Order Differential Equations In this chapter we will consider systems of n first-order linear differential equations in n unknown functions:

x1

t = a11
tx1
t + a12
tx2
t + · · · + a1n
txn
t + g1
t x2

t = a21
tx1
t + a22
tx2
t + · · · + a2n
txn
t + g2
t   xn

t = an1
tx1
t + an2
tx2
t + · · · + ann
txn
t + gn
t 361

362

CHAPTER 10

Systems of Linear Differential Equations

Let

⎛ ⎜ ⎜ A
t = ⎜ ⎝ ⎛ ⎜ ⎜ X
t = ⎜ ⎝

a11
t a21
t  

a12
t a22
t  

an1
t

an2
t

x1
t x2
t  

⎞ ⎟ ⎟ ⎟ ⎠

··· ···

a1n
t a2n
t  

··· · · · ann
t ⎛

⎜ ⎜ and G
t = ⎜ ⎝

xn
t

⎞ ⎟ ⎟ ⎟ ⎠ g1
t g2
t  

⎞ ⎟ ⎟ ⎟ ⎠

gn
t

Differentiate a matrix by differentiating each element, so ⎛
⎞ x1
t ⎜ x2

t ⎟ ⎜ ⎟
X
t = ⎜  ⎟  ⎝  ⎠ xn

t Matrix differentiation follows the “normal” rules we learn in calculus. For example,
X
tY
t
= X

tY
t + X
tY

t in which the order of the factors must be maintained. Now the system of differential equations is X

t = A
tX
t + G
t

(10.1)

or X
= AX + G This system is nonhomogeneous if G
t = O for at least some t, in which O denotes the n × 1 zero matrix ⎛ ⎞ 0 ⎜ 0 ⎟ ⎜ ⎟ ⎜  ⎟  ⎝  ⎠ 0 If G
t = O for all the relevant values of t, then the system is homogeneous, and we write just X
= AX A solution of X
= AX + G is any n × 1 matrix of functions that satisfies this matrix equation.

EXAMPLE 10.1

The 2 × 2 system x1
= 3x1 + 3x2 + 8 x2
= x1 + 5x2 + 4e3t

10.1 Theory of Systems of Linear First-Order Differential Equations can be written



x1 x2




 =

One solution is

 X
t =

3 1

3 5



x1 x2



 +

8 4e3t

3e2t + e6t − 4e3t − 103 −e2t + e6t + 23

363

 

 

as can be verified by substitution into the system. In terms of individual components, this solution is x1
t = 3e2t + e6t − 4e3t −

10 3

2 x2
t = −e2t + e6t +  3

On February 20, 1962, John Glenn became the first American to orbit the Earth. His flight lasted nearly five hours and included three complete circuits of the globe. This and subsequent Mercury orbitings paved the way for the space shuttle program, which now includes shuttles launched from the NASA Kennedy Space Center to carry out experiments under zero gravity, as well as delivery of personnel and equipment to the developing international space station. Ultimate goals of space shuttle missions include studying how humans function in a zero-gravity environment over extended periods of time, scientific observations of phenomena in space and on our own planet, and the commercial development of space. Computation of orbits and forces involved in shuttle missions involves the solution of large systems of differential equations.

364

CHAPTER 10

Systems of Linear Differential Equations

Initial conditions for the system (10.1) have the form ⎛ ⎞ x1
t0  ⎜ x2
t0  ⎟ ⎜ ⎟ X
t0  = ⎜ ⎟ = X0   ⎝ ⎠  xn
t0  in which X0 is a given n × 1 matrix of constants. The initial value problem we will consider for systems is the problem: X
= AX + G

X
t0  = X0 

(10.2)

This is analogous to the initial value problem x
= ax + g

x
t0  = x0

for single first-order equations. Theorem 1.3 gave criteria for existence and uniqueness of solutions of this initial value problem. The analogous result for the initial value problem (10.2) is given by the following. THEOREM 10.1

Existence and Uniqueness

Let I be an open interval containing t0 . Suppose each aij
t and gj
t are continuous on I. Let X0 be a given n × 1 matrix of real numbers. Then the initial value problem X
= AX + G

X
t0  = X0

has a unique solution defined for all t in I.

EXAMPLE 10.2

Consider the initial value problem x1
= x1 + tx2 + cos
t x2
= t3 x1 − et x2 + 1 − t x1
0 = 2 x2
0 = −5 This is the system

 X
=

1 t3

t −et

with



 X0 =



cos
t 1−t

X+

2 −5

 

 

This initial value problem has a unique solution defined for all real t, because each aij
t and gj
t are continuous for all real t. We will now determine what we must look for to find all solutions of X
= AX + G. This will involve a program that closely parallels that for the single first-order equation x
= ax + g, beginning with the homogeneous case.

10.1 Theory of Systems of Linear First-Order Differential Equations

10.1.1

365

Theory of the Homogeneous System X
= AX

We begin with the homogeneous system X
= AX. Because solutions of X
= AX are n × 1 matrices of real functions, these solutions have an algebraic structure, and we can form linear combinations (finite sums of scalar multiples of solutions). In the homogeneous case, any linear combination of solutions is again a solution.

THEOREM 10.2

Let 1      k be solutions of X
= AX, all defined on some open interval I. Let c1      ck be any real numbers. Then the linear combination c1 1 + · · · + ck k is also a solution of X
= AX, defined on I. Proof

Compute
c1 1 + · · · + ck k 
= c1 1
+ · · · + ck k
= c1 A1 + · · · + ck Ak = A
c1 1 + · · · + ck k 

Because of this, the set of all solutions of X
= AX has the structure of a vector space, called the solution space of this system. It is not necessary to have a background in vector spaces to follow the discussion of solutions of X
= AX that we are about to develop. However, for those who do have this background we will make occasional reference to show how ideas fit into this algebraic framework. In a linear combination c1 1 + · · · + ck k of solutions, any j that is already a linear combination of the other solutions is unnecessary. For example, suppose 1 = a2 2 + · · · + ak k . Then c1 1 + · · · + ck k = c1
a2 2 + · · · + ak k  + c2 2 + · · · + ck k =
c1 a2 + c2 2 + · · · +
c1 ak + ck k  In this case any linear combination of 1  2      k is actually a linear combination of just 2      k , and 1 is not needed. 1 is redundant in the sense that, if we have 2      k , then we have 1 also. We describe this situation by saying that the functions 1  2      k are linearly dependent. If no one of the functions is a linear combination of the others, then these functions are called linearly independent.

DEFINITION 10.1 Linear Dependence

Solutions 1  2      k , of X
= AX, defined on an interval I, are linearly dependent on I if one solution is a linear combination of the others on this interval. Linear Independence

Solutions 1  2      k , of X
= AX, defined on an interval I, are linearly independent on I if no solution in this list is a linear combination of the others on this interval.

Thus a set of solutions is linearly independent if it is not linearly dependent. Linear dependence of functions is a stronger condition than linear dependence of vectors. For vectors in Rn , V1 is a linear combination of V2 and V3 if V1 = aV2 + bV3 for some real

366

CHAPTER 10

Systems of Linear Differential Equations

numbers a and b. In this case V1  V2  V3 are linearly dependent. But for solutions 1  2  3 of X
= AX, 1 is a linear combination of 2 and 3 if there are numbers a and b such that 1
t = a2
t + b2
t for all t in the relevant interval, perhaps the entire real line. It is not enough to have this condition hold for just some values of t.

EXAMPLE 10.3

Consider the system




X = It is routine to check that



1
t =

−2e3t e3t

1 1

−4 5

 X



 and 2
t =


1 − 2te3t te3t



are solutions, defined for all real values of t. These solutions are linearly independent on the entire real line, since neither is a constant multiple of the other (for all real t). The function  
11 − 6te3t 3
t =
−4 + 3te3t is also a solution. However, 1  2  3 are linearly dependent, because, for all real t, 3
t = −41
t + 32
t This means that 3 is a linear combination of 1 and 2 , and the list of solutions 1  2  3 , although longer, carries no more information about the solution of X
= AX than the list of solutions 1  2 . If  is a solution of X
= AX, then  is an n × 1 column matrix of functions: ⎛ ⎞ f1
t ⎜ f2
t ⎟ ⎜ ⎟ 
t = ⎜  ⎟  ⎝  ⎠ fn
t For any choice of t, say t = t0 , this is an n × 1 matrix of real numbers which can be thought of as a vector in Rn . This point of view, and some facts about determinants, provides us with a test for linear independence of solutions of X
= AX. The following theorem reduces the question of linear independence of n solutions of X
= AX, to a question of whether an n × n determinant of real numbers is nonzero.

THEOREM 10.3

Test for Linear Independence of Solutions

Suppose that

⎛

⎜ ⎜ 1
t = ⎜ ⎝

11
t 21
t   n1
t

⎞

⎛

⎜ ⎟ ⎜ ⎟ ⎟  2
t = ⎜ ⎝ ⎠

12
t 22
t   n2
t

⎞

⎛

⎜ ⎟ ⎜ ⎟ ⎟      n
t = ⎜ ⎝ ⎠

1n
t 2n
t   nn
t

⎞ ⎟ ⎟ ⎟ ⎠

10.1 Theory of Systems of Linear First-Order Differential Equations

367

are solutions of X
= AX on an open interval I. Let t0 be any number in I. Then 1. 1  2      n are linearly independent on I if and only if 1
t0      n
t0  are linearly independent, when considered as vectors in Rn . 2. 1  2      n are linearly independent on I if and only if    11
t0  12
t0  · · · 1n
t0      21
t0  22
t0  · · · 2n
t0       = 0        · · ·     n1
t0  n2
t0  · · · nn
t0   Conclusion (2) is an effective test for linear independence of n solutions of X
= AX on an open interval. Evaluate each solution at some point of the interval. Each j
t0  is an n × 1 (constant) column matrix. Evaluate the determinant of the n × n matrix having these columns. If this determinant is nonzero, then the solutions are linearly independent; if it is zero, they are linearly dependent. Another way of looking at (2) of this theorem is that it reduces a question of linear independence of n solutions of X
= AX, to a question of linear independence of n vectors in Rn . This is because the determinant in (2) is nonzero exactly when its row (or column) vectors are linearly independent.

EXAMPLE 10.4

From the preceding example,   −2e3t 1
t = e3t are solutions of

 and

2
t =


1 − 2te3t te3t





 1 −4 X = X 1 5 on the entire real line, which is an open interval. Evaluate these solutions at some convenient point, say t = 0:     −2 1 1
0 = and 2
0 =  1 0


Use these as columns of a 2 × 2 matrix and evaluate its determinant:    −2 1     1 0  = −1 = 0 Therefore 1 and 2 are linearly independent solutions. A proof of Theorem 10.3 makes use of the uniqueness of solutions of the initial value problem (Theorem 10.1). For (1), let t0 be any point in I. Suppose first that 1      n are linearly dependent on I. Then one of the solutions is a linear combination of the others. By reordering if necessary, say 1 is a linear combination of 2      n . Then there are numbers c2      cn so that

Proof

1
t = c2 2
t + · · · + cn n
t

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CHAPTER 10

Systems of Linear Differential Equations

for all t in I. In particular, 1
t0  = c2 2
t0  + · · · + cn n
t0  This implies that the vectors 1
t0      n
t0  are linearly dependent vectors in Rn . Conversely, suppose that 1
t0      n
t0  are linearly dependent in Rn . Then one of these vectors is a linear combination of the others. Again, as a convenience, suppose 1
t0  is a linear combination of 2
t0      n
t0 . Then there are numbers c2      cn such that 1
t0  = c2 2
t0  + · · · + cn n
t0  Define 
t = 1
t − c2 2
t − · · · − cn n
t for all t in I. Then  is a linear combination of solutions of X
= AX, hence is a solution. Further ⎛ ⎞ 0 ⎜ 0 ⎟ ⎜ ⎟ 
t0  = ⎜  ⎟  ⎝  ⎠ 0 Therefore, on I,  is a solution of the initial value problem X
= AX But the zero function

X
t0  = O ⎛

⎜ ⎜ 
t = ⎜ ⎝

0 0  

⎞ ⎟ ⎟ ⎟ ⎠

0 is also a solution of this initial value problem. Since this initial value problem has a unique solution, then for all t in I, ⎛ ⎞ 0 ⎜ 0 ⎟ ⎜ ⎟ 
t = 
t = ⎜  ⎟  ⎝  ⎠ 0 Therefore

⎛ ⎜ ⎜ 
t = 1
t − c2 2
t − · · · − cn n
t = ⎜ ⎝

0 0  

⎞ ⎟ ⎟ ⎟ ⎠

0 for all t in I, which means that 1
t = c2 2
t + · · · + cn n
t for all t in I. Therefore 1 is a linear combination of 2      n , hence 1  2      n are linearly dependent on I.

10.1 Theory of Systems of Linear First-Order Differential Equations

369

Conclusion (2) follows from (1) and the fact that n vectors in Rn are linearly independent if and only if the determinant of the n × n matrix having these vectors as columns is nonzero. Thus far we know how to test n solutions of X
= AX for linear independence, if A is n × n. We will now show that n linearly independent solutions are enough to determine all solutions of X
= AX on an open interval I. We saw a result like this previously when it was found that two linear independent solutions of y

+ p
xy
+ q
xy = 0 determine all solutions of this equation. THEOREM 10.4

Let A = aij
t be an n × n matrix of functions that are continuous on an open interval I. Then 1. The system X
= AX has n linearly independent solutions defined on I. 2. Given any n linearly independent solutions 1      n defined on I, every solution on I is a linear combination of 1      n . By (2), every solution of X
= AX, defined on I , must be of the form c1 1 + c2 2 + · · · + cn n . For this reason, this linear combination, with 1      n any n linearly independent solutions, is called the general solution of X
= AX on I. To prove that there are matrices ⎛ 1 ⎜ 0 ⎜ ⎜ E
1 = ⎜ 0 ⎜  ⎝ 

Proof

n linearly independent solutions, define the n × 1 constant ⎞

⎛

⎟ ⎜ ⎟ ⎜ ⎟
2 ⎜ ⎟E = ⎜ ⎟ ⎜ ⎠ ⎝

0

0 1 0  

⎞

⎛

⎟ ⎜ ⎟ ⎜ ⎟ ⎜
n ⎟E = ⎜ ⎟ ⎜ ⎠ ⎝

0

0 0 0  

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

1

Pick any t0 in I. We know from Theorem 10.1 that the initial value problem X
= AX

X
t0  = E
j

has a unique solution j defined on I, for j = 1 2    n. These solutions are linearly independent by Theorem 10.3, because the way the initial conditions were chosen, the n × n matrix whose columns are these solutions evaluated at t0 is In , with determinant 1. This proves (1). To prove (2), suppose now that 1      n are any n linearly independent solutions of X
= AX, defined on I. Let be any solution. We want to prove that is a linear combination of 1      n . Pick any t0 in I. We will first show that there are numbers c1      cn such that
t0  = c1 1
t0  + · · · + cn n
t0  Now
t0 , and each j
t0 , is an n × 1 column matrix of constants. Form the n × n matrix S using 1
t0      n
t0  as its columns, and consider the system of n linear algebraic equations in n unknowns ⎛ ⎞ c1 ⎜ c2 ⎟ ⎟ ⎜ S ⎜  ⎟ =
t0  (10.3) ⎝  ⎠ cn

370

CHAPTER 10

Systems of Linear Differential Equations

The columns of S are linearly independent vectors in Rn , because 1      n are linearly independent. Therefore S is nonsingular, and the system (10.3) has a unique solution. This solution gives constants c1      cn such that
t0  = c1 1
t0  + · · · + cn n
t0  We now claim that
t = c1 1
t + · · · + cn n
t for all t in I. But observe that and c1 1 + · · · + cn n are both solutions of the initial value problem X
= AX

X
t0  =
t0 

Since this problem has a unique solution, then
t = c1 1
t + · · · + cn n
t for all t in I , and the proof is complete. In the language of linear algebra, the solution space of X
= AX has dimension n, the order of the coefficient matrix A. Any n linearly independent solutions form a basis for this vector space.

EXAMPLE 10.5

Previously we saw that  1
t =

−2e3t e3t



 and 2
t =


1 − 2te3t te3t



are linearly independent solutions of 


X =

1 1

−4 5

 X

Because A is 2 × 2 and we have 2 linearly independent solutions, the general solution of this system is     −2e3t
1 − 2te3t 
t = c1 + c  2 e3t te3t The expression on the right contains every solution of this system. In terms of components, x1
t = −2c1 e3t + c2
1 − 2te3t  x2
t = c1 e3t + c2 te3t  We will now make a useful observation. In the last example, form a 2 × 2 matrix
t having 1
t and 2
t as columns:   −2e3t
1 − 2te3t 


t = e3t te3t

10.1 Theory of Systems of Linear First-Order Differential Equations  Now observe that, if C =

c1 c2

371

 , then 


tC = =

  −2e3t
1 − 2te3t c1 e3t te3t c2 

  c1 −2e3t + c2
1 − 2te3t

c1 e3t  + c2 te3t      −2e3t
1 − 2te3t + c = c1 2 e3t te3t = c1 1
t + c2 2
t

The point is that we can write the general solution c1 1 + c2 2 compactly as
tC, with


t a square matrix having the independent solutions as columns, and C a column matrix of arbitrary constants. We call a matrix formed in this way a fundamental matrix for the system X
= AX. In terms of this fundamental matrix, the general solution is X
t =
tC. We can see that
tC also satisfies the matrix differential equation X
= AX. Recall that we differentiate a matrix by differentiating each element of the matrix. Then, because C is a constant matrix,   −6e3t e3t − 6te3t

tC
=
t
C = C 3e3t e3t + 3te3t Now compute A

tC =
A
tC    1 −4 −2e3t
1 − 2te3t = C 1 5 e3t te3t   −2e3t − 4e3t
1 − 2te3t − 4te3t C = −2e3t + 5e3t
1 − 2te3t + 5te3t   −6e3t
1 − 6te3t =  3e3t e3t + 3te3t Therefore

tC
= A

tC as occurs if
tC is a solution of X
= AX.

DEFINITION 10.2

is a fundamental matrix for the n × n system X
= AX if the columns of are linearly independent solutions of this system.

Writing the general solution of X
= AX as X
t = C is particularly convenient for solving initial value problems.

372

CHAPTER 10

Systems of Linear Differential Equations

EXAMPLE 10.6

Solve the initial value problem






X =

−4 5

1 1



 

X
0 =

−2 3

 

We know from Example 10.5 that the general solution if X
t =
tC, where   −2e3t
1 − 2te3t


t =  e3t te3t We need to choose C so that

 X
0 =
0C =

−2 3

 

Putting t = 0 into , we must solve the algebraic system     −2 1 −2 C=  1 0 3 The solution is

 C=  =

−2 1 0 1

−1   1 −2 0 3     1 −2 3 =  2 3 4

The unique solution of the initial value problem is therefore   3 
t =
t 4   −2e−3t − 8te3t =  3e3t + 4te3t In this example
0−1 could be found by linear algebra methods (Sections 7.8.1 and 8.7) or by using a software package.

10.1.2 General Solution of the Nonhomogeneous System X
= AX+ G Solutions of the nonhomogeneous system X
= AX + G do not have the algebraic structure of a vector space, because linear combinations of solutions are not solutions. However, we will show that the general solution of this system (an expression containing all possible solutions) is the sum of the general solution of the homogeneous system X
= AX, and any particular solution of the nonhomogeneous system. This is completely analogous to Theorem 2.5 for the second order equation y

+ p
xy
+ q
xy = f
x.

THEOREM 10.5

Let be a fundamental matrix for X
= AX, and let p be any solution of X
= AX + G. Then the general solution of X
= AX + G is X = C +  p , in which C is an n × 1 matrix of arbitrary constants.

10.1 Theory of Systems of Linear First-Order Differential Equations Proof

373

First, C +  p is a solution of the nonhomogeneous system, because
C +  p 
=
C
+ p
= A
C + A p + G = A
C +  p  + G

Now let  be any solution of X
= AX + G. We claim that  −  p is a solution of X
= AX. To see this, calculate
 −  p 
= 
− p
= A + G −
A p + G = A
 −  p  Since C is the general solution of X
= AX, there is a constant n × 1 matrix K such that  −  p = K. Then  = K +  p , completing the proof. We now know what to look for in solving a system of n linear, first-order differential equations in n unknown functions. For the homogeneous system, X
= AX, we look for n linearly independent solutions to form a fundamental matrix
t. For the nonhomogeneous system X
= AX + G, we first find the general solution C of X
= AX, and any particular solution p of X
= AX + G. The general solution of X
= AX + G is then C +  p . This is an overall strategy. Now we need ways of implementing it and actually producing fundamental matrices and particular solutions for given systems.

SECTION 10.1

PROBLEMS

In each of Problems 1 through 5, (a) verify that the given functions satisfy the system, (b) form a fundamental matrix
t for the system, (c) write the general solution in the form
tC, carry out this product, and verify that the rows of
tC are the components of the given solution, and (d) find the unique solution satisfying the initial conditions. 1.

4.

 x1
t = 2e

x2
t = c1 e

x1
= 5x1 + 3x2  x2
= x1 + 3x2

x2
t = 2c1 e4t cos
t − sin
t + 2c2 e4t cos
t + sin
t

  √  15t 15t + c2 sin  2 2

3t/2

√   √  √ 15t 15t − cos + 15 sin 2 2  √  √  √ 15t 15t + 15 cos  2 2

x1
0 = −2 x2
0 = 7

x1
t = c1 e4t cos
t + c2 e4t sin
t 5.

x1
= 5x1 − 4x2 + 4x3  x2
= 12x1 − 11x2 + 12x3  x3
= 4x1 − 4x2 + 5x3 

x1
0 = −2 x2
0 = 1 x1

t = 3x1 + 8x2  x2

t = x1 − x2 x1
t = 4c1 e

c1 cos



x1
= 2x1 + x2  x2
= −3x1 + 6x2

√
1+2 3t

√

− c2 e3t/2 sin

x1
0 = 0 x2
0 = 4

3.

3t/2



x1
t = −c1 e2t + 3c2 e6t  x2
t = c1 e2t + c2 e6t

2.

x1
= x1 − x2  x2
= 4x1 + 2x2

√
1−2 3t

x1
t = c1 et + c3 e−3t  x2
t = c2 e2t + c3 e−3t 

+ 4c2 e  √ √ √ √ x2
t =
−1 + 3c1 e
1+2 3t +
−1 − 3c2 e
1−2 3t

x3
t =
c2 − c1 et + c3 e−3t 

x1
0 = 2 x2
0 = 2

x1
0 = 1 x2
0 = −3 x3
0 = 5

CHAPTER 10

374

10.2

Systems of Linear Differential Equations

Solution of X
= AX when A is Constant Consider the system X
= AX, with A an n × n matrix of real numbers. In the case y
= ay, with a constant, we get exponential solutions y = ceax . This suggests we try a similar solution for the system. Try X = et , with  an n × 1 matrix of constants to be determined, and  a number to be determined. Substitute this proposed solution into the differential equation to get et = A
et  This requires that A =  We should therefore choose  as an eigenvalue of A, and  as an associated eigenvector. We will summarize this discussion.

THEOREM 10.6

Let A be an n × n matrix of real numbers. Then et is a nontrivial solution of X
= AX if and only if  is an eigenvalue of A, with associated eigenvector . We need n linearly independent solutions to form a fundamental matrix. We will have these if we can find n linearly independent eigenvectors, whether or not some eigenvalues may be repeated. THEOREM 10.7

Let A be an n × n matrix of real numbers. Suppose A has eigenvalues 1      n , and suppose there are associated eigenvectors 1      n that are linearly independent. Then 1 e1 t      n en t are linearly independent solutions of X
= AX, on the entire real line. We know that each j ej t is a nontrivial solution. The question is whether these solutions are linearly independent. Form the n × n matrix having these solutions, evaluated at t = 0, as its columns. This matrix has n linearly independent columns 1      n , and therefore has a nonzero determinant. By Theorem 10.3(2), 1 e1 t      n en t are linearly independent on the real line.

Proof

EXAMPLE 10.7

Consider the system

 X
=

4 3

2 3

 X 

 1 . These 1 − 23 eigenvectors are linearly independent (originating from distinct eigenvalues), so we have two linearly independent solutions,     1 1 t e and e6t  1 − 23 A has eigenvalues 1 and 6, with corresponding eigenvectors

1





and

10.2 Solution of X
= AX when A is Constant We can write the general solution as



X
t = c1

1



 et + c2

− 23

1 1

375

 e6t 

Equivalently, we can write the fundamental matrix   et e6t


t =  − 23 et e6t In terms of , the general solution is X
t =
tC. In terms of components, x1
t = c1 et + c2 e6t 3 x2
t = − c1 et + c2 e6t  2

EXAMPLE 10.8

Solve the system

⎛

5 X
= ⎝ 12 4

−4 −11 −4

⎞ 4 12 ⎠ X 5

The eigenvalues of A are −3 1 1. Even though one eigenvalue is repeated, A has three linearly independent eigenvectors. They are: ⎛ ⎞ 1 ⎝ 3 ⎠ associated with eigenvalue − 3 1 and

⎛

⎞ 1 ⎝ 1 ⎠ 0

⎛ and

⎞ −1 ⎝ 0 ⎠ associated with 1 1

This gives us three linearly independent solutions ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 1 1 −1 ⎝ 3 ⎠ e−3t  ⎝ 1 ⎠ et  ⎝ 0 ⎠ et  1 0 1 A fundamental matrix is

⎛

e−3t ⎝ 3e−3t


t = e−3t

et et 0

⎞ −et 0 ⎠ et

The general solution is X
t =
tC.

EXAMPLE 10.9 A Mixing Problem

Two tanks are connected by a series of pipes, as shown in Figure 10.1. Tank 1 initially contains 20 liters of water in which 150 grams of chlorine are dissolved. Tank 2 initially contains 50 grams of chlorine dissolved in 10 liters of water.

376

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Systems of Linear Differential Equations Pure water: 3 liters/min

Mixture: 3 liters/min

Tank 1

Mixture: 2 liters/min

Tank 2

Mixture: 4 liters/min

Mixture: 1 liter/min

FIGURE 10.1

Beginning at time t = 0, pure water is pumped into tank 1 at a rate of 3 liters per minute, while chlorine/water solutions are interchanged between the tanks and also flow out of both tanks at the rates shown. The problem is to determine the amount of chlorine in each tank at any time t > 0. At the given rates of input and discharge of solutions, the amount of solution in each tank will remain constant. Therefore, the ratio of chlorine to chlorine/water solution in each tank should, in the long run, approach that of the input, which is pure water. We will use this observation as a check of the analysis we are about to do. Let xj
t be the number of grams of chlorine in tank j at time t. Reading from Figure 10.1, rate of change of xj
t = xj

t = rate in minus rate out      gram liter x  gram liter ·0 +3 · 2 =3 min liter min 10 liter      liter liter x1 gram x  gram −2 −4 · · 1 min 20 liter min 20 liter =−

6 3 x1 + x2  20 10

Similarly, with the dimensions excluded, x2

t = 4

x x 4 x1 4 − 3 2 − 2 = x1 − x2  20 10 10 20 10

The system is X
= AX, with  A=

− 103 1 5

3 10 − 25

 

The initial conditions are x1
0 = 150 x2
0 = 50 or  X
0 =

150 50

 

10.2 Solution of X
= AX when A is Constant

377

The eigenvalues of A are − 101 and − 35 , and corresponding eigenvectors are, respectively,  3    −1 2 and  1 1 These are linearly independent and we can write the fundamental matrix   3 e−t/10 −e−3t/5 2 


t = e−t/10 e−3t/5 The general solution is X
t =
tC. To solve the initial value problem, we must find C so that  3    −1 150 2 X
0 = =
0C = C 50 1 1 Then

 C=  =

3 2

−1

1

1

2 5 2 −5

2 5 3 5

−1  

150 50

150 50





 =

80 −30

 

The solution of the initial value problem is  3   e−t/10 −e−3t/5 80 2 X
t = −30 e−t/10 e−3t/5   120e−t/10 + 30e−3t/5 =  80e−t/10 − 30e−3t/5 Notice that x1
t → 0 and x2
t → 0 as t → , as we expected.

10.2.1

Solution of X
= AX when A has Complex Eigenvalues

Consider a system X
= AX. If A is a real matrix, the characteristic polynomial of A has real coefficients. It may, however, have some complex roots. Suppose  =  + i is a complex eigenvalue, with eigenvector . Then A = , so A =  But A = A if A has real elements, so A =  This means that  =  − i is also an eigenvalue, with eigenvector . This means that et and et can be used as two of the n linearly independent solutions needed to form a fundamental matrix. This resulting fundamental matrix will contain some complex entries. There is nothing wrong with this. However, sometimes it is convenient to have a fundamental matrix with only real entries. We will show how to replace these two columns, involving complex numbers, with two other linearly independent solutions involving only real quantities. This can be done for any pair of columns arising from a pair of complex conjugate eigenvalues.

378

CHAPTER 10

Systems of Linear Differential Equations

THEOREM 10.8

Let A be an n × n real matrix. Let  + i be a complex eigenvalue with corresponding eigenvector U + iV, in which U and V are real n × 1 matrices. Then et U cos
t − V sin
t and et U sin
t + V cos
t are real linearly independent solutions of X
= AX.

EXAMPLE 10.10

Solve the system X
= AX, with ⎛

2 A=⎝ 0 0

0 −2 2

⎞ 1 −2 ⎠  0

√ √ The eigenvalues are 2 −1 + 3i −1 − 3i. Corresponding eigenvectors are, respectively, ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 1√ 1 1 √ ⎝ 0 ⎠  ⎝ −2 3i ⎠  ⎝ 2 3i ⎠  √ √ 0 −3 + 3i −3 − 3i One solution is

⎛

⎞ 1 ⎝ 0 ⎠ e2t 0

and two other solutions are ⎛ ⎞ 1√ √ ⎝ −2 3i ⎠ e
−1+ 3it √ −3 + 3i

⎛ and

⎝

1 √

⎞

⎠ e
−1− 2 3i √ −3 − 3i

√ 3tt



These three solutions are linearly independent and can be used as columns of a fundamental matrix ⎞ ⎛ √ √ e
−1+ 3it√ e
−1− 3tt√ e2t √ √ ⎟ ⎜
−1− 3tt

1
t = ⎝ 0 −2 √ 3ie
−1+ 3it 2 3ie ⎠ √ √ √
−1+ 3it
−1− 3tt 0
−3 + 3ie
−3 − 3ie However, we can also produce a real fundamental matrix as follows. First write ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ 0√ 1√ 1 ⎝ −2 3i ⎠ = ⎝ 0 ⎠ + i ⎝ −2 3 ⎠ = U + iV √ √ −3 3 −3 + 3i

10.2 Solution of X
= AX when A is Constant with

⎛

⎞ 1 U=⎝ 0 ⎠ −3

Then

379

⎛

and

⎞ 0√ ⎠ V = ⎝ −2 √ 3  3

⎞ 1√  √ √ √  ⎝ −2 3i ⎠ e
−1+ 3it =
U + iV e−t cos
3t + ie−t sin
3t √ −3 + 3i  √ √ √ = Ue−t cos
3t − Ve−t sin
3t + i Ve−t cos
3t √  (10.4) +Ue−t sin
3t  ⎛

And ⎛

⎞ 1  √ √ √  √ ⎝ 2 3i ⎠ e
−1− 3tt =
U − iV e−t cos
3t − ie−t sin
3t √ −3 − 3i  √ √ √ = Ue−t cos
3t − Ve−t sin
3t − i Ve−t cos
3t √  + Ue−t sin
3t  (10.5)

The functions (10.4) and (10.5) are solutions, so any linear combination of these is also a solution. Taking their sum and dividing by 2 yields the solution √ √ 1
t = Ue−t cos
3t − Ve−t sin
3t And taking their difference and dividing by 2i yields the solution √ √ 2
t = Ve−t cos
3t + Ue−t sin
3t Using these, together with the solution found from the eigenvalue 2, we can form the fundamental matrix √ √ ⎛ 2t ⎞ e−t cos
3t e−t sin
3t e √ √ √ √ −t −t ⎠

2
t = ⎝ 0 2 3e √ sin
√ 3t √ √−2 3e√ cos
3t √ −t −t 0 e −3 cos
3t − 3 sin
3t e 3 cos
3t − 3 sin
3t Either fundamental matrix can be used to write the general solution, X
t = 1
tC or X
t =

2
tK. However, the latter involves only real numbers and real-valued functions. A proof of the theorem follows the reasoning of the example, and is left to the student.

10.2.2

Solution of X
= AX when A does not have n Linearly Independent Eigenvectors

We know how to produce a fundamental matrix for X
= AX when A has n linearly independent eigenvectors. This certainly occurs if A has n distinct eigenvalues, and may even occur when A has repeated eigenvalues. However, we may encounter a matrix A having repeated eigenvalues, for which there are not n linearly independent eigenvectors. In this case we cannot yet write a fundamental matrix. This section is devoted to a procedure to follow in this case to find a fundamental matrix. We will begin with two examples and then make some general remarks.

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Systems of Linear Differential Equations

EXAMPLE 10.11

We will solve the system X
= AX, with



A=

1 −3

3 7

  

A has one eigenvalue 4 of multiplicity 2. Eigenvectors all have the form 

1 1

 , with  = 0.

A does not have two linearly independent eigenvectors. We can immediately write one solution   1 1
t = e4t  1   1 We need another solution. Write E1 = and attempt a second solution 1 2
t = E1 te4t + E2 e4t  in which E2 is a 2 × 1 constant matrix to be determined. For this to be a solution, we need to have 2

t = A2
t: E1 e4t + 4te4t  + 4E2 e4t = AE1 te4t + AE2 e4t  Divide this equation by e4t to get E1 + 4E1 t + 4E2 = AE1 t + AE2  But AE1 = 4E1 , so the terms having t as a factor cancel and we are left with AE2 − 4E2 = E1  Write this equation as  If E2 =

a b


A − 4I2 E2 = E1 



, this is the linear system of two equations in two unknowns: 
A − 4I2 

or



−3 −3

3 3 



a b





1 1

=

a b 



 =

 

1 1

 

s This system has general solution E2 = , in which s can be any nonzero number. Let 1+3s 3   1 s = 1 to get E2 = and hence the second solution 4 3

 2
t = E1 te + E2 e = 4t

 =

1+t 4 +t 3

4t

 e4t 

1 1



 te + 4t

1 4 3

 e4t

10.2 Solution of X
= AX when A is Constant

381

If we use 1
0 and 2
0 as columns to form the matrix   1 1  1 43 then this matrix has determinant 1/3, hence 1 and 2 are linearly independent by Theorem 10.3(2). Therefore 1
t and 2
t are linearly independent and can be used as columns of a fundamental matrix   4t
1 + te4t e 


t = e4t
43 + te4t The general solution of X
= AX is X
t =
tC. The procedure followed in this example is similar in spirit to solving the differential equation y

− 5y
+ 6y = e3x by undetermined coefficients. We are tempted to try yp
x = ae3x , but this will not work because e3x is a solution of y

− 5y
+ 6y = 0. We therefore try yp
x = axe3x , multiplying the first attempt ae3x by x. The analogous step for the system was to try the second solution 2
t = E1 te4t + E2 e4t . We will continue to explore the case of repeated eigenvalues with another example.

EXAMPLE 10.12

Consider the system X
= AX, in which ⎛

⎞ −2 −1 −5 0 ⎠ A = ⎝ 25 −7 0 1 3

A has eigenvalue scalar ⎛ −2⎞with multiplicity 3, and corresponding eigenvectors are ⎛ all nonzero ⎞ −1 −1 multiples of ⎝ −5 ⎠. This gives us one solution of X
= AX. Denoting ⎝ −5 ⎠ = E1 , we 1 1 have one solution ⎞ ⎛ −1 1
t = ⎝ −5 ⎠ e−2t = E1 e−2t  1 We need three linearly independent solutions. We will try a second solution of the form 2
t = E1 te−2t + E2 e−2t  in which E2 is a 3 × 1 matrix to be determined. Substitute this proposed solution into X
= AX to get E1 e−2t − 2te−2t  + E2 −2e−2t  = AE1 te−2t + AE2 e−2t  Upon dividing by the common factor of e−2t , and recalling that AE1 = −2E1 , this equation becomes E1 − 2tE1 − 2E2 = −2tE1 + AE2  or AE2 + 2E2 = E1 

382

CHAPTER 10

Systems of Linear Differential Equations

We can write this equation as
A + 2I3 E2 = E1  or

⎛

0 ⎝ 25 0

−1 −5 1

⎛ ⎞ ⎞ −5 −1 0 ⎠ E2 = ⎝ −5 ⎠  5 1

⎞  With E2 = ⎝  ⎠, this is the nonhomogeneous system ! ⎛ ⎞⎛ ⎞ ⎛ ⎞ 0 −1 −5  −1 ⎝ 25 −5 0 ⎠ ⎝  ⎠ = ⎝ −5 ⎠  0 1 5 ! 1 ⎛ ⎞ −s with general solution ⎝ 1 − 5s ⎠, in which s can be any number. For a specific solution, s choose s = 1 and let ⎛ ⎞ −1 E2 = ⎝ −4 ⎠  1 ⎛

This gives us the second solution 2
t = E1 te−2t + E2 e−2t ⎛ ⎞ ⎛ ⎞ −1 −1 = ⎝ −5 ⎠ te−2t + ⎝ −4 ⎠ e−2t 1 1 ⎛ ⎞ −1 − t = ⎝ −4 − 5t ⎠ e−2t  1+t We need one more solution. Try for a solution of the form 1 3
t = E1 t2 e−2t + E2 te−2t + E3 e−2t  2 We want to solve for E3 . Substitute this proposed solution into X
= AX to get 1 E1 te−2t − t2 e−2t  + E2 e−2t − 2te−2t  + E3 −2e−2t  = AE1 t2 e−2t + AE2 te−2t + AE3 e−2t  2 ⎛ ⎞ 1 Divide this equation by e−2t and use the fact that AE1 = −2E1 and AE2 = ⎝ 3 ⎠ to get −1 ⎞ ⎛ 1 E1 t − E1 t2 + E2 − 2E2 t − 2E3 = −E1 t2 + ⎝ 3 ⎠ t + AE3  (10.6) −1

10.2 Solution of X
= AX when A is Constant Now

383

⎡⎛

⎞ ⎛ ⎞⎤ ⎛ ⎞ −1 −1 1 E1 t − 2E2 t = ⎣⎝ −5 ⎠ − 2 ⎝ −4 ⎠⎦ t = ⎝ 3 ⎠ t 1 1 −1

so equation (10.6) has three terms cancel, and it reduces to E2 − 2E3 = AE3  Write this equation as
A + 2I3 E3 = E2 or

⎛

0 ⎝ 25 0

⎛ ⎞ ⎞ −1 −5 −1 −5 0 ⎠ E3 = ⎝ −4 ⎠  1 5 1

with general solution

⎛

1−25s 25

⎞

⎟ ⎜ ⎝ 1 − 5s ⎠ s in which s can be any number. Choosing s = 1, we ⎛ 24 − 25 ⎜ E3 = ⎝ −4 1 A third solution is

can let ⎞ ⎟ ⎠

⎛ 24 ⎞ ⎞ ⎛ ⎞ − 25 −1 −1 1⎝ ⎟ ⎜ 2 −2t −2t ⎠ ⎝ ⎠ −5 t e + −4 te + ⎝ −4 ⎠ e−2t 3
t = 2 1 1 1 ⎞ ⎛ 24 − 25 − t − 21 t2 ⎟ ⎜ = ⎝ −4 − 4t − 25 t2 ⎠ e−2t  ⎛

1 + t + 21 t2 To show that 1 , 2 and 3 are linearly independent, Theorem 10.3(2) is convenient. Form the 3 × 3 matrix having these solutions, evaluated at t = 0, as columns: ⎛ ⎞ −1 −1 − 24 25 ⎜ ⎟ ⎝ −5 −4 −4 ⎠  1 1 1 The determinant of this matrix is − 251 , so this matrix is nonsingular and the solutions are linearly independent. We can use these solutions as columns of a fundamental matrix  24  ⎛ ⎞ − 25 − t − 21 t2 e−2t
−1 − te−2t −e−2t   ⎟ ⎜


t = ⎝ −5e−2t
−4 − 5te−2t −4 − 4t − 25 t2 e−2t ⎠    1 + t + 21 t2 e−2t e−2t
1 + te−2t The general solution of X
= AX is X
t =
tC.

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CHAPTER 10

Systems of Linear Differential Equations

These examples suggest a procedure which we will now outline in general. Begin with a system X
= AX, with A an n × n matrix of real numbers. We want the general solution, so we need n linearly independent solutions. Case 1—A has n linearly independent eigenvectors. Use these eigenvectors to write n linearly independent solutions and use these as columns of a fundamental matrix. (This case may occur even if A does not have n distinct eigenvalues). Case 2—A does not have n linearly independent eigenvectors. Let the eigenvalues of A be 1      n At least one eigenvalue must be repeated, because if A has n distinct eigenvectors, the corresponding eigenvectors must be linearly independent, putting us back in case 1. Suppose 1      r are the distinct eigenvalues, while r+1      n repeat some of these first r eigenvalues. If Vj is an eigenvector corresponding to j for j = 1    r, we can immediately write r linearly independent solutions 1
t = V1 e1 t      r
t = Vr er t  Now work with the repeated eigenvalues. Suppose  is a repeated eigenvalue, say  = 1 with multiplicity k. We already have one solution corresponding to , namely 1 . To be consistent in notation with the examples just done, denote V1 = E1 and 1 = 1 . Then 1
t = V1 e1 t = E1 et is one solution corresponding to . For a second solution corresponding to , let 2
t = E1 tet + E2 et  Substitute this proposed solution into X
= AX and solve for E2 . If k = 2, this yields a second solution corresponding to  and we move on to another multiple eigenvalue. If k ≥ 3, we do not yet have all the solutions corresponding to , so we attempt 1 3
t = E1 t2 et + E2 tet + E3 et  2 Substitute 3
t into the differential equation and solve for E3 to get a third solution corresponding to . If  ≥ 4, continue with 1 1 E t3 et + E2
tt2 et + E3 tet + E4 et  3! 1 2! substitute into the differential equation and solve for E4 , and so on. Eventually, we reach 4
t =

k
t =

1 1 E1 tk−1 et + E tk−2 et + · · · + Ek−1 tet + Ek et 
k − 1!
k − 2! 2

substitute into the differential equation and solve for Ek . This procedure gives, for an eigenvalue  of multiplicity k, k linearly independent solutions of X
= AX. Repeat the procedure for each eigenvalue until n linearly independent solutions have been found.

10.2.3 Solution of X
= AX by Diagonalizing A We now take a different tack and attempt to exploit diagonalization. Consider the system ⎞ ⎛ −2 0 0 0 ⎠ X X
= ⎝ 0 4 0 0 −6

10.2 Solution of X
= AX when A is Constant

385

The constant coefficient matrix A is a diagonal matrix, and this system really consists of three independent differential equations, each involving just one of the variables: x1
= −2x1  x2
= 4x2  x3
= −6x3  Such a system is said to be uncoupled. Each equation is easily solved independent of the others, obtaining x1 = c1 e−2t  x2 = c2 e4t  x3 = c3 e−6t  The system is uncoupled because the coefficient matrix A is diagonal. Because of this, we can immediately write the eigenvalues −2 4 −6 of A, and find the corresponding eigenvectors ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 1 0 0 ⎝ 0 ⎠⎝ 1 ⎠⎝ 0 ⎠ 0 0 1 Therefore X
= AX has fundamental matrix ⎛ −2t e


t = ⎝ 0 0

0 e4t 0

0 0

⎞ ⎠

e−6t

and the general solution is X
t =
tC. However we wish to approach this system, the point is that it is easy to solve because A is a diagonal matrix. Now in general A need not be diagonal. However, A may be diagonalizable (Section 9.2). This will occur exactly when A has n linearly independent eigenvectors. In this event, we can form a matrix P, whose columns are eigenvectors of A, such that ⎛ ⎞ 1 0 · · · 0 ⎜ 0 2 · · · 0 ⎟ ⎜ ⎟ P−1 AP = D = ⎜    ⎟ ⎝   · · ·  ⎠ 0

0

···

n

D is the diagonal matrix having the eigenvalues 1      n down its main diagonal. This will hold even if some of the eigenvalues have multiplicities greater than 1, provided that A has n linearly independent eigenvectors. Now make the change of variables X = PZ in the differential equation X
= AX. First compute X
=
PZ
= PZ
= AX = APZ so Z
= P−1 APZ = DZ The uncoupled system Z
= DZ can be solved by inspection. A fundamental matrix for Z
= DZ is ⎞ ⎛ t e1 0 ··· ··· 0 ⎜ 0 e2 t · · · · · · 0 ⎟ ⎟ ⎜ ⎜     ⎟  

D
t = ⎜  ⎟ · · ·  ⎜ ⎟ ⎝ 0 0 ⎠ 0 · · · en−1 t 0 0 ··· 0 en t

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CHAPTER 10

Systems of Linear Differential Equations

and the general solution of Z
= DZ is Z
t = D
tC Then X
t = PZ
t = P D
tC is the general solution of the original system X
= AX. That is,
t = P D
t is a fundamental matrix for X
= AX. In this process we need P, whose columns are eigenvectors of A, but we never actually need to calculate P−1 .

EXAMPLE 10.13

Solve

 X
=

3 1



3 5

X

The eigenvalues and associated eigenvectors of A are     −3 1 2, and 6  1 1 Because A has distinct eigenvalues, A is diagonalizable. Make the change of variables X = PZ, where   −3 1 P=  1 1 This transforms X
= AX to Z
= DZ, where  2 D= 0



0 6



This uncoupled system has fundamental matrix  2t e

D
t = 0 Then X
= AX has fundamental matrix




t = P D
t =  =

−3 1

0 e6t

1 1

−3e2t e2t





e6t e6t



e2t 0  

0 e6t



The general solution of X
= AX is X
t =
tC.

10.2.4 Exponential Matrix Solutions of X
= AX A first-order differential equation y
= ay has general solution y
x = ceax . At the risk of stretching the analogy too far, we might conjecture whether there might be a solution eAt C for a matrix differential equation X
= AX. We will now show how to define the exponential matrix eAt to make sense out of this conjecture. For this section, let A be an n × n matrix of real numbers.

10.2 Solution of X
= AX when A is Constant

387

The Taylor expansion of the real exponential function, et = 1 + t +

1 2 1 3 t + t +···  2! 3!

suggests the following.

DEFINITION 10.3

Exponential Matrix

The exponential matrix eAt is the n × n matrix defined by eAt = In + At +

1 2 2 1 3 3 A t + A t +···  2! 3!

It can be shown that this series converges for all real t, in the sense that the infinite series of elements in the i j place converges. Care must be taken in computing with exponential matrices, because matrix multiplication is not commutative. The analogue of the relationship eat ebt = e
a+bt is given by the following. THEOREM 10.9

Let B be an n × n real matrix. Suppose AB = BA. Then e
A+Bt = eAt eBt  Because A is a constant matrix,

d 1 2 2 1 3 3 1 4 4 d At e = I + At + A t + A t + A t + · · · dt dt n 2! 3! 4! 1 1 = A + A2 t + A3 t2 + A4 t3 + · · · 2! 3!

1 2 2 1 3 3 1 4 4 = A In + At + A t + A t + A t + · · · 2! 3! 4! = AeAt  The derivative of eAt , obtained by differentiating each element, is the product AeAt of two n × n matrices, and has the same form as the derivative of the scalar exponential function eat . One ramification of this derivative formula is that, for any n × 1 constant matrix K, eAt K is a solution of X
= AX. LEMMA 10.1

For any real n × 1 constant matrix K, eAt K is a solution of X
= AX. Proof

Compute 

t =
eAt K
= AeAt K = A
t

Even more, eAt is a fundamental matrix for X
= AX.

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THEOREM 10.10

eAt is a fundamental matrix for X
= AX. Proof

Let Ej be the n × 1 matrix with 1 in the j 1 place and all other entries zero: ⎛ ⎞ 0 ⎜ 0 ⎟ ⎜ ⎟ ⎜  ⎟ ⎜  ⎟ ⎜ ⎟ ⎟ Ej = ⎜ ⎜ 1 ⎟ ⎜  ⎟ ⎜  ⎟ ⎜ ⎟ ⎝ 0 ⎠ 0

Then eAt Ej is the j th column of eAt . This column is a solution of X
= AX by the lemma. Further, the columns of eAt are linearly independent by Theorem 10.3(2), because if we put t = 0, we get eAt = In , which has a nonzero determinant. Thus is a fundamental matrix for X
= AX. In theory, then, we can find the general solution eAt C of X
= AX, if we can compute eAt . This, however, can be a daunting task. As an example, for an apparently simple matrix such as   1 2 A=  −2 4 we find using a software package that eAt = √ √ √ √  5t/2  √4 e5t/2 sin
7t/2 e cos
7t/2 + 7e5t/2 sin
7t/2 7  √ √ √ √ e5t/2 cos
7t/2 − 7e5t/2 sin
7t/2 − √47 e5t/2 sin
7t/2 This is a fundamental matrix for X
= AX. It would be at least as easy, for this A, to find the  √ 4√ eigenvalues of A, which are 25 ± 21 i 7, then find corresponding eigenvectors, , 3±i 7 and use these to obtain a fundamental matrix. We will now pursue an interesting line of thought. We claim that, even though eAt may be tedious or even impractical to compute for a given A, it is often possible to compute the product eAt K, for carefully chosen K, as a finite sum, and hence generate solutions to X
= AX. To do this we need the following. LEMMA 10.2

Let A be an n × n real matrix and K an n × 1 real matrix. Let  be any number. Then 1. eIn t K = et K. 2. eAt K = et e
A−In t K. Proof

For (1), since
In m = In for any positive integer m, we have

1 1 In t 2 2 3 3 e K = In + In t +
In  t +
In  t + · · · K 2! 3!

1 2 2 1 3 3 = 1 + t +  t +  t + · · · In K = et K 2! 3!

10.2 Solution of X
= AX when A is Constant

389

For (2), first observe that In and A − In commute, since In
A − In  = 
In A − 
In 2  = 
A − In  =
A − In 
In  Then, using Theorem 10.9, eAt K = eAt+In t−In t K = e
A−In t eIn t K = e
A−In t et K = et e
A−In t K Now suppose we want to solve X
= AX. Let 1      r be the distinct eigenvalues of A, and let j have multiplicity mj . Then m1 + · · · + mr = n For each j , find as many linearly independent eigenvectors as possible. For j , this can be any number from 1 to mj inclusive. If this yields n linearly independent eigenvectors, then we can write the general solution as a sum of eigenvectors j times exponential functions ej t , and we do not need eAt . Thus suppose some j has multiplicity mj ≥ 2, but there are fewer than mj linearly independent eigenvectors. Find an n × 1 constant matrix K1 that is linearly independent from the eigenvectors found for j and such that
A − j In K1 = O,

but


A − j In 2 K1 = O

Then eAt K1 is a solution of X
= AX. Further, because of the way K1 was chosen

 eAt K1 = ej t e
A−j In t K1 = ej t K1 +
A − j In K1 t  with all other terms of the series for e
A−j In t K1 vanishing because
A − j In 2 K1 = O forces
A − j In m K1 = O for m ≥ 2. We can therefore compute ej t K1 as a sum of just two terms. If we now have mj solutions corresponding to j , then leave this eigenvalue and move on to any others that do not yet have as many linearly independent solutions as their multiplicity. If we do not yet have mj solutions corresponding to j , then find a constant n × 1 matrix K2 such that
A − j In K2 = O

and
A − j In 2 K2 = O

but
A − j In 3 K2 = O Then eAt K2 is a solution of X
= AX, and we can compute this solution as a sum of just three terms:

1 eAt K2 = ej t e
A−j In t K2 = ej t K2 +
A − j In K2 t +
A − j In 2 t2  2! The other terms in the infinite series for eAt K2 vanish because
A − j In 3 K2 =
A − j In 4 K2 = · · · = O If this gives us mj solutions associated with j , move on to another eigenvalue for which we do not yet have as many solutions as the multiplicity of the eigenvalue. If not, produce an n × 1 constant matrix K3 such that
A − j In K3 = O,
A − j In 2 K3 = O and
A − j In 3 K3 = O but
A − j In 4 K3 = 0 Then eAt K3 can be computed as a sum of four terms.

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Keep repeating this process. For j it must terminate after at most mj − 1 steps, because we began with at least one eigenvector associated with j and then produced more solutions to obtain a total of mj linearly independent solutions associated with j . Once these are obtained, we move on to another eigenvalue for which we have fewer solutions than its multiplicity, and repeat this process for that eigenvalue, and so on. Eventually we generate a total of n linearly independent solutions, thus obtaining the general solution of X
= AX.

EXAMPLE 10.14

Consider X
= AX, where

⎛

2 ⎜ 0 A=⎜ ⎝ 0 0

1 2 0 0

0 1 2 0

⎞ 3 1 ⎟ ⎟ 4 ⎠ 4 ⎛

⎞ 9 ⎜ 6 ⎟ ⎟ The eigenvalues are 4 2 2 2. Associated with 4 we find the eigenvector ⎜ ⎝ 8 ⎠, so one 4 solution of X
= AX is ⎛ ⎞ 9 ⎜ 6 ⎟ 4t ⎟ 1
t = ⎜ ⎝ 8 ⎠e  4 Associated with 2 we find that every eigenvector has the form ⎛ ⎞  ⎜ 0 ⎟ ⎜ ⎟ ⎝ 0 ⎠ 0 A second solution is

⎛

⎞ 1 ⎜ 0 ⎟ 2t ⎟ 2
t = ⎜ ⎝ 0 ⎠e  0

Now find a 4 × 1 constant matrix K1 such that
A − 2I4 K1 = O, but
A − 2I4 2 K1 = O. First compute ⎛ ⎞2 0 1 0 3 ⎜ 0 0 1 1 ⎟ ⎟
A − 2I4 2 = ⎜ ⎝ 0 0 0 4 ⎠ 0 0 0 2 ⎛ ⎞ 0 0 1 7 ⎜ 0 0 0 6 ⎟ ⎟ =⎜ ⎝ 0 0 0 8 ⎠ 0 0 0 4

10.2 Solution of X
= AX when A is Constant

391

Solve
A − 2I4 2 K1 = O to find solutions of the form ⎛ ⎞  ⎜  ⎟ ⎜ ⎟ ⎝ 0 ⎠ 0 We will choose the solution

⎛

⎞ 0 ⎜ 1 ⎟ ⎟ K1 = ⎜ ⎝ 0 ⎠ 0

to avoid duplicating the eigenvector already found associated with 2. Then ⎛ ⎞⎛ ⎞ ⎛ ⎞ 0 1 0 3 0 1 ⎜ 0 0 1 1 ⎟⎜ 1 ⎟ ⎜ 0 ⎟ ⎟⎜ ⎟ ⎜ ⎟
A − 2I4 K1 = ⎜ ⎝ 0 0 0 4 ⎠ ⎝ 0 ⎠ = ⎝ 0 ⎠ = O 0 0 0 2 0 0 as required. Thus form the third solution 3
t = eAt K1 = e2t K1 +
A − 2I4 K1 t ⎡⎛ ⎞ ⎛ ⎞⎛ ⎞ ⎤ 0 0 1 0 3 0 ⎢⎜ 1 ⎟ ⎜ 0 0 1 1 ⎟ ⎜ 1 ⎟ ⎥ ⎜ ⎟ ⎜ ⎟⎜ ⎟ ⎥ = e2t ⎢ ⎣⎝ 0 ⎠ + ⎝ 0 0 0 4 ⎠ ⎝ 0 ⎠ t⎦ 0 0 0 0 2 0 ⎡⎛ ⎞ ⎛ ⎞ ⎤ ⎛ ⎞ 0 1 t ⎢⎜ ⎟ ⎜ 0 ⎟ ⎥ ⎜ 1 ⎟ 2t 2t ⎢⎜ 1 ⎟ ⎟ t⎥ = ⎜ ⎟e  +⎜ = e ⎣⎝ 0 ⎠ ⎝ 0 ⎠ ⎦ ⎝ 0 ⎠ 0 0 0 The three solutions found up to this point are linearly independent. Now we need a fourth solution. It must come from the eigenvalue 2, because 4 has multiplicity 1 and we have one solution corresponding to this eigenvalue. Look for K2 such that
A − 2I4 K2 = O

and
A − 2I4 2 K2 = O

but
A − 2I4 3 K2 = O First compute

⎛

0 ⎜ 0 3 ⎜
A − 2I4  = ⎝ 0 0 Solutions of
A − 2I4 3 K2 = O are of the form ⎛  ⎜  ⎜ ⎝ ! 0

0 0 0 0 ⎞ ⎟ ⎟ ⎠

0 0 0 0

⎞ 18 12 ⎟ ⎟ 16 ⎠ 8

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We will choose

⎛

⎞ 1 ⎜ 1 ⎟ ⎟ K2 = ⎜ ⎝ 1 ⎠ 0

to avoid duplicating previous choices. Of course other choices are possible. It is verify that
A − 2I4 K2 = O and
A − 2I4 2 K2 = O. Thus form the fourth solution

1 At 2t 2 2 4
t = e K2 = e K2 +
A − 2In K2 t +
A − 2In  t 2! ⎞ ⎞⎛ ⎞ ⎛ ⎡⎛ ⎛ ⎞⎛ 1 0 1 0 3 1 0 0 1 7 1 ⎟ ⎜ 0 0 1 1 ⎟⎜ 1 ⎟ ⎢⎜ ⎜ 0 0 0 6 ⎟⎜ 1 1 2t ⎢⎜ 1 ⎟ ⎟t+ ⎜ ⎟⎜ ⎟⎜ +⎜ = e ⎣⎝ 1 ⎠ ⎝ 0 0 0 4 ⎠⎝ 1 ⎠ 2 ⎝ 0 0 0 8 ⎠⎝ 1 0 0 0 0 2 0 0 0 0 4 0 ⎛ 1 2 ⎞ 1+t + 2t ⎜ ⎟ 2t 1+t ⎟e  =⎜ ⎝ ⎠ 1

routine to

⎞ ⎤ ⎟ 2⎥ ⎟t ⎥ ⎠ ⎦

0 We now have four linearly independent solutions, hence the general solution. We can also write the fundamental matrix ⎛ 4t ⎞ 9e e2t te2t
1 + t + 21 t2 e2t ⎜ 6e4t 0 e2t ⎟
1 + te2t ⎟


t = ⎜ ⎝ 8e4t 0 ⎠ 2t 0 e 4e4t

SECTION 10.2

0

0

0

PROBLEMS

In each of Problems 1 through 5, find a fundamental matrix for the system and use it to write the general solution. These coefficients matrices for these systems have real, distinct eigenvalues. 1. x1
= 3x1  x2
= 5x1 − 4x2 2. x1
= 4x1 + 2x2  x2
= 3x1 + 3x2 3. x1
= x1 + x2  x2
= x1 + x2 4. x1
= 2x1 + x2 − 2x3  x2
= 3x1 − 2x2  x3
= 3x1 − x2 − 3x3 5. x1
= x1 +2x2 +x3  x2
= 6x1 −x2  x3
= −x1 −2x2 −x3 In each of Problems 6 through 11, find a fundamental matrix for the system and use it to solve the initial value problem. The matrices of these systems have real, distinct eigenvalues.

6. x1
= 3x1 − 4x2  x2
= 2x1 − 3x2  x1
0 = 7 x2
0 = 5 7. x1
= x1 − 2x2  x2
= −6x1  x1
0 = 1 x2
0 = −19 8. x1
= 2x1 − 10x2  x2
= −x1 − x2  x1
0 = −3 x2
0 =6 9. x1
= 3x1 − x2 + x3  x2
= x1 + x2 − x3  x3
= x1 − x2 + x3  x1
0 = 1 x2
0 = 5 x3
0 = 1 10. x1
= 2x1 + x2 − 2x3  x2
= 3x1 − 2x2  x3
= 3x1 + x2 − 3x3  x1
0 = 1 x2
0 = 7 x3
0 = 3 11. x1
= 2x1 + 3x2 + 3x3  x2
= −x2 − 3x3  x3
= 2x3  x1
0 = 9 x2
0 = −1 x3
0 = −3

10.2 Solution of X
= AX when A is Constant 12. Show that the change of variables z = ln
t for t > 0 transforms the system tx1
= ax1 + bx2  tx2
= cx1 + dx2 into a linear system X
= AX, assuming that a b c and d are real constants. 13. Use the idea of Problem 12 to solve the system tx1
= 6x1 + 2x2  tx2
= 4x1 + 4x2  14. Solve the system tx1
= −x1 − 3x2  tx2
= x1 − 5x2  In each of Problems 15 through 19, find a real-valued fundamental matrix for the system X
= AX, with A the given matrix.   2 −4 15. 1 2   0 5 16. −1 −2   3 −5 17. 1 −1 ⎞ ⎛ 1 −1 1 0 ⎠ 18. ⎝ 1 −1 1 0 −1 ⎞ ⎛ −2 1 0 0 ⎠ 19. ⎝ −5 0 0 3 −2 In each of Problems 20 through 23, find a real-valued fundamental matrix for the system X
= AX, with A the given matrix. Use this to solve the initial value problem, with X
0 the given n × 1 matrix.     2 3 2  20. −5 1 8     1 3 −2  21. 10 5 −3     5 2 −5  22. 0 1 −2 ⎞ ⎛ ⎛ ⎞ 3 −3 1 7 23. ⎝ 2 −1 0 ⎠  ⎝ 4 ⎠ 1 −1 1 3 24. Can a matrix with at least one complex, non-real element have only real eigenvalues? If not, give a proof. It it can, give an example. In each of Problems 25 through 30, find a fundamental matrix for the system X
= AX, using the method of Section 10.2.2, with A the given matrix.

 25.  26.

3 0

2 3

2 5

0 2

⎛

2 27. ⎝ 0 0 ⎛ 1 28. ⎝ 0 4 ⎛ 1 ⎜ 0 29. ⎜ ⎝ 0 0 ⎛ ⎜ 30. ⎜ ⎝

0 0 0 −1

393

  ⎞ 6 9 ⎠ 2 ⎞

5 8 −1 5 1 8

0 0 ⎠ 1

5 3 3 0

−2 0 0 0 1 0 0 −2

0 1 0 0

⎞ 6 4 ⎟ ⎟ 4 ⎠ 1 ⎞ 0 0 ⎟ ⎟ 1 ⎠ 0

In each of Problems 31 through 35, find the general solution of the system X
= AX, with A the given matrix, and use this general solution to solve the initial value problem, for the given n × 1 matrix X
0. Use the method of Section 10.2.2 for these problems.     7 −1 5 31.  1 5 3     2 0 4  32. 3 5 2 ⎞ ⎞ ⎛ ⎛ 0 −4 1 1 33. ⎝ 0 2 −5 ⎠  ⎝ 4 ⎠ 12 0 0 −4 ⎞ ⎞ ⎛ ⎛ 2 −5 2 1 3 ⎠  ⎝ −3 ⎠ 34. ⎝ 0 −5 4 0 0 −5 ⎞ ⎛ ⎛ ⎞ 1 −2 0 0 2 ⎟ ⎜ 1 −1 0 ⎜ 0 ⎟ ⎜ −2 ⎟ ⎟ 35. ⎜  ⎝ 0 0 5 −3 ⎠ ⎝ 1 ⎠ 0 0 3 −1 4 In each of Problems 36 through 40, find the general solution of the system by diagonalizing the coefficient matrix. 36. x1
= −2x1 + x2  x2
= −4x1 + 3x2 37. x1
= 3x1 + 3x2  x1
= x1 + 5x2 38. x1
= x1 + x2  x2
= x1 + x2 39. x1
= 6x1 + 5x2  x2
= x1 + 2x2 40. x1
= 3x1 − 2x2  x2
= 9x1 − 3x2

CHAPTER 10

394

Systems of Linear Differential Equations

In each of Problems 41-45, solve the system X
= AX, with A the matrix of the indicated problem, by finding eAt . 41. 42. 43. 44. 45.

A A A A A

10.3

as as as as as

in in in in in

Problem Problem Problem Problem Problem

In each of Problems 46-50, solve the initial value problem of the referred problem, using the exponential matrix.

25. 26. 27. 28. 29.

46. 47. 48. 49. 50.

Problem Problem Problem Problem Problem

31. 32. 33. 34. 35.

Solution of X
= AX + G We now turn to the nonhomogeneous system X

t = A
tX
t + G
t, assuming that the elements of the n × n matrix A
t, and the n × 1 matrix G
t, are continuous on some interval I, which may be the entire real line. Recall that the general solution of X
= AX + G has the form X
t =
tC + p
t, where


t is an n × n fundamental matrix for the homogeneous system X
= AX, C is an n × 1 matrix of arbitrary constants, and p is a particular solution of X
= AX + G. At least when A is a real, constant matrix, we have a strategy for finding . We will concentrate in this section on strategies for finding a particular solution p .

10.3.1 Variation of Parameters Recall the variation of parameters method for second-order differential equations. If y1
x and y2
x form a fundamental set of solutions for y


x + p
xy

x + q
xy
x = 0 then the general solution of this homogeneous equation is yh
x = c1 y1
x + c2 y2
x To find a particular solution yp
x of the nonhomogeneous equation y


x + p
xy

x + q
xy
x = f
x replace the constants in yh by functions and attempt to choose u
x and v
x so that yp
x = u
xy1
x + v
xy2
x is a solution. The variation of parameters method for the matrix equation X
= AX + G follows the same idea. Suppose we can find a fundamental matrix for the homogeneous system X
= AX. The general solution of this homogeneous system is then Xh
t =
tC, in which C is an n × 1 matrix of arbitrary constants. Look for a particular solution of X
= AX + G of the form p
t =
tU
t in which U
t is an n × 1 matrix of functions of t which is to be determined. Substitute this proposed solution into the differential equation to get
U
= A
U + G, or


U + U
=
A U + G

10.3 Solution of X
= AX + G

395

Now is a fundamental matrix for X
= AX, so
= A . Therefore
U =
A U and the last equation reduces to

U
= G Since is a fundamental matrix, the columns of are linearly independent. This means that is nonsingular, so the last equation can be solved for U
to get U
= −1 G As in the case of second order differential equations, we now have the derivative of the function we want. Then
U
t = −1
tG
tdt in which we integrate a matrix by integrating each element of the matrix. Once we find a suitable U
t, we have a particular solution p
t =
tU
t of X
= AX + G. The general solution of this nonhomogeneous equation is then X
t =
tC +
tU
t in which C is an n × 1 matrix of constants.

EXAMPLE 10.15

Solve the system




X =

1 −10 −1 4



 X+

et sin
t

 


First we need a fundamental matrix for of A are −1 and 6, with  X= AX.  The eigenvalues  5 −2 associated eigenvectors, respectively, and . Therefore a fundamental matrix for 1 1
X = AX is   −t −2e6t 5e 


t = e−t e6t

We find (details provided at the end of the example) that   1 2et et 

−1
t = 7 −e−6t 5e−6t Compute U

t = −1
tG
t =

1 7

1 = 7 Then

 

et −e−6t

2et 5e−6t



e2t + 2et sin
t −e−5t + 5e−6t sin
t

et sin
t  



 t  2t  sin
tdt  e−5tdt + 2  e −6t − e dt + 5 e sin
tdt   1 2t e + 17 et sin
t − cos
t 14 =  1 −5t 5 −6t e + 259 e −6 sin
t − cos
t 35

U
t =




−1
tG
tdt =

1 7



396

CHAPTER 10

Systems of Linear Differential Equations

The general solution of X
= AX + G is



5e−t e−t

X
t =
tC +
tU
t =  +  =

5e−t e−t

5e−t e−t

−2e6t e6t −2e6t e6t



 C

1 2t e + 17 et sin
t − cos
t 14 1 −5t 5 −6t e + 259 e −6 sin
t − cos
t 35





−2e6t e6t

3 t e 10 1 t e 10

C+

+ 35 sin
t − 25 cos
t 37 37





+ 371 sin
t − 376 cos
t



If we want to write the solution in terms of the component functions, let C =



c1 c2

 to obtain

25 3 t 35 e + sin
t − cos
t 10 37 37 6 1 1 x2
t = c1 e−t + c2 e6t + et + sin
t − cos
t 10 37 37 x1
t = 5c1 e−t − 2c2 e6t +

Although the coefficient matrix A in this example was constant, this is not a requirement of the variation of parameters method. In the example we needed −1
t. Standard software packages will produce this inverse. We could also proceed as follows, reducing
t and recording the row operations beginning with the identity matrix, as discussed in Section 7.8.1: ⎞ ⎛  −t 6t −2e ⎠ ⎝ 1 0  5e  −t 0 1  e e6t add − 15 (row 1) to row 2

⎛ 1

0

− 15

1

⎝ multiply row 1 by 15 et

⎛ ⎜ ⎝

multiply row 2 by 57 e−6t

add 25 e7t (row 2) to row 1

1 t e 5

0

− 15

1

   

1 t e 5

0

− 17 e−6t

5 −6t e 7

0

⎛ ⎜ ⎝

⎞ − 25 e7t ⎟ ⎠ 7 6t e 5

1

⎛ ⎜ ⎝

⎞ −2e6t ⎠ 7 6t e 5

  5e−t   0

1 t e 7

2 t e 7

− 17 e−6t

5 −6t e 7

   

⎞ − 25 e7t ⎟ ⎠ 1

1 0

   

1 0

⎞ 0 ⎟ ⎠ 1

Since the last two columns for I2 , the first two columns are −1
t.

10.3 Solution of X
= AX + G

397

Variation of Parameters and the Laplace Transform There is a connection between the variation of parameters method and the Laplace transform. Suppose we want a particular solution p of X
= AX + G, in which A is an n × n real matrix. The variation of parameters method is to find a particular solution p
t =
tU
t , where
t is a fundamental matrix for X
= AX. Explicitly,
U
t = −1
tG
tdt We can choose a particular U
t by carrying out this integration from 0 to t:
t U
t =

−1
sG
sds 0

Then p
t =
t




t 0

−1
sG
sds =




t 0


t −1
sG
sds

In this equation can be any fundamental matrix for X
= AX. In particular, suppose we choose
t = eAt . This is sometimes called the transition matrix for X
= AX, since it is a fundamental matrix such that
0 = In . Now −1
s = e−As , so


t −1
s = eAt e−As = eA
t−s =
t − s and p
t =




t 0


t − sG
sds

This equation has the same form as the Laplace transform convolution of and G, except that in the current setting these are matrix functions. Now define the Laplace transform of a matrix to be the matrix obtained by taking the Laplace transform of each of its elements. This extended Laplace transform has many of the same computational properties as the Laplace transform for scalar functions. In particular, we can define the convolution integral
t


t − sG
sds


t ∗ G
t = 0

In terms of this convolution, p
t =
t ∗ G
t This is a general formula for a particular solution of X
= AX + G when
t = eAt .

EXAMPLE 10.16

Consider the system  X
= We find that

 e = At

1 1

−4 5


1 − 2te3t te3t



 X+

e2t t

−4te3t
1 + 2te3t

 

 =
t

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CHAPTER 10

Systems of Linear Differential Equations

A particular solution of X
= AX + G is given by
t p
t =


t − sG
sds 0

  2s  −4
t − se3
t−s e dt
1 + 2
t − se3
t−s s 0 
t
1 − 2t + 2s e3t e−s − 4
t − s e3t se−3s = ds
t − s e3t e−s +
1 + 2t − 2s e3t se−3s 0  t  3t −s 3t −3s
t


1 − 2t + 2s e e − 4 − s e se ds 0 t =


t − s e3t e−s +
1 + 2t − 2s e3t se−3s ds 0   −3e2t + 89 e3t − 229 e3t t − 49 t − 278 27 =  e2t + 119 e3t t − 28 e3t − 19 t + 271 27

=


t


1 − 2
t − se3
t−s
t − se3
t−s

The general solution of X
= AX + G is  
1 − 2te3t −4te3t X
t = C te3t
1 + 2te3t   −3e2t + 89 e3t − 229 e3t t − 49 t − 278 27 +  e2t + 119 e3t t − 28 e3t − 19 t + 271 27

10.3.2 Solution of X
= AX+ G by Diagonalizing A Consider the case that A is a constant, diagonalizable matrix. Then A has n linearly independent eigenvectors. These form columns of a nonsingular matrix P such that ⎛ ⎞ 1 0 · · · 0 ⎜ 0 2 · · · 0 ⎟ ⎜ ⎟ P−1 AP = D = ⎜    ⎟  ⎝   ···  ⎠ 0 0 · · · n with eigenvalues down the main diagonal in the order corresponding to the eigenvector columns of P. As we did in the homogeneous case X
= AX, make the change of variables X = PZ. Then the system X
= AX + G becomes X
= PZ
= A
PZ + G or PZ
=
APZ + G Multiply this equation on the left by P−1 to get Z
=
P−1 APZ + P−1 G = DZ + P−1 G This is an uncoupled system of the form z
1 = 1 z1 + f1
t z
2 = 2 z2 + f2
t   z
n = n zn + fn
t

10.3 Solution of X
= AX + G where

⎛

⎞

f1
t f2
t  

⎜ ⎜ P−1 G
t = ⎜ ⎝

399

⎟ ⎟ ⎟ ⎠

fn
t ⎛ ⎜ ⎜ Solve these n first-order differential equations independently, form Z
t = ⎜ ⎝

z1
t z2
t  

⎞ ⎟ ⎟ ⎟, and then ⎠

zn
t the solution of X
= AX + G is X
t = PZ
t. Unlike diagonalization in solving the homogeneous system X
= AX, in this nonhomogeneous case we must explicitly calculate P−1 in order to determine P−1 G
t, the nonhomogeneous term in the transformed system.

EXAMPLE 10.17

Consider




X =

3 1



3 5



8 4e3t

X+

  

The eigenvalues of A are 2 and 6, with eigenvectors, respectively, Let

 P=

−3 1

Then

 −1

P AP =



1 1

2 0





0 6



Compute  P

−1

=

− 41

1 4 3 4

1 4

 

The transformation X = PZ transforms the original system into     2 0 8 Z + P−1 Z
= 0 6 4e3t      − 41 41 2 0 8 = Z+  1 3 0 6 4e3t 4

or




Z =

2 0

0 6



 Z+

4

−2 + e3t 2 + 3e3t

 

−3 1



 and

1 1

 .

400

CHAPTER 10

Systems of Linear Differential Equations

This is the uncoupled system z
1 = 2z1 − 2 + e3t z
2 = 6z2 + 2 + 3e3t  Solve these linear first-order differential equations independently: z1
t = c1 e2t + e3t + 1 1 z2
t = c2 e6t − e3t −  3 Then

 Z
t =

Then

 X
t = PZ
t =  =

−3 1

c1 e2t + e3t + 1 c2 e6t − e3t − 13

1 1



 

c1 e2t + e3t + 1 c2 e6t − e3t − 13

−3c1 e2t + c2 e6t − 4e3t − 103

 

c1 e2t + c2 e6t + 23     −4e3t − 103 −3e2t e6t = C+  2 e2t e6t

(10.7)

3

This is the general solution of X
= AX + G. It is the general solution of X
= AX, plus a particular solution of X
= AX + G. Indeed,   −3e2t e6t e2t e6t is a fundamental matrix for the homogeneous system X
= AX. To illustrate the solution of an initial value problem, suppose we want the solution of       2 3 3 8  X
0 =  X+ X
= −7 1 5 4e3t Since we have the general solution (10.7) of this system, all we need to do is determine C so that this initial condition is satisfied. We need     −4 − 103 −3 1 X
0 = C+ 2 1 1 3  22      −3 −3 1 2 = C+ =  2 1 1 −7 3

This is the equation  PC =

2 + 223 −7 − 23



 =

28 3 − 233

 

10.3 Solution of X
= AX + G We already have P−1 , so C=P

 −1

28 3 − 233



 =  =

− 41 1 4

− 174



1 4 3 4



28 3 − 233

401





41 − 12

The initial value problem has the unique solution     17   −4e3t − 103 −4 −3e2t e6t
X = + 41 2 e2t e6t − 12 3  51  41 6t e2t − 12 e − 4e3t − 103 4 =  41 6t − 174 e2t − 12 e + 23

SECTION 10.3

PROBLEMS

In each of Problems 1 through 5, use variation of parameters to find the general solution of X
= AX + G, with A and G as given.     −3et 5 2  1. −2 1 e3t     1 2 −4  2. 3t 1 −2     2e6t 7 −1  3. 6t 1 5 6te ⎞ ⎞ ⎛ 2t ⎛ e cos
3t 2 0 0 ⎠ −2 4. ⎝ 0 6 −4 ⎠  ⎝ −2 0 4 −2 ⎞ ⎛ ⎛ ⎞ 0 1 0 0 0 t ⎜ 4 3 0 0 ⎟ ⎜ −2e ⎟ ⎟ ⎜ ⎟ 5. ⎜ ⎝ 0 0 3 0 ⎠⎝ 0 ⎠ −1 2 9 1 et In each of Problems 6 through 9, use variation of parameters to solve the initial value problem X
= AX + G X
0 = X0 , with A, G and X0 as given (in that order).       0 2 0 2   6. 3 10t 5 2       5 −4 2et −1 7.   t 4 −3 3 2e ⎞ ⎞ ⎛ ⎞ ⎛ ⎛ 2t 10e 5 2 −3 1 2 4 ⎠  ⎝ 6e2t ⎠  ⎝ 11 ⎠ 8. ⎝ 0 −2 0 0 1 −e2t

⎛

1 9. ⎝ 3 4

−3 −5 7

⎞ ⎛ ⎞ ⎞ ⎛ te−2t 0 6 −2t ⎠⎝ 2 ⎠ 0 ⎠  ⎝ te −2 3 t2 e−2t

10. Recall that a transition matrix for X
= AX is a fundamental matrix 
t such that 
0 = In . (a) Prove that, for a transition matrix
t, −1
t =


−t and
t + s =
t
s for real s and t. (b) Suppose
t is any fundamental matrix for X
= AX. Prove that 
t =
t −1
0 is a transition matrix. That is, 
t is a fundamental matrix and 
0 = In . In each of Problems 11, 12 and 13, verify that the given matrix is a fundamental matrix for the system and use it to find a transition matrix.   2et e6t 11. x1
= 4x1 + 2x2  x2
= 3x1 + 3x2  −3et e6t   −5t 2e
1 + 5te−5t 5
12. x1 = −10x2  x2 = x1 −10x2  5 −5t te e−5t 2 2 13. x1
= 5x1 − 4x2 + 4x3  x2
= 12x1 − 11x2 + 12x3 , ⎞ ⎛ −3t et 0 e
0 et ⎠ x3 = 4x1 − 4x2 + 5x3  ⎝ 3e−3t e−3t −et et In each of Problems 14 through 18, find the general solution of the system by diagonalization. The general solutions of the associated homogeneous systems X
= AX were requested in Problems 36 through 40 of Section 10.2.

CHAPTER 10

402

Systems of Linear Differential Equations

14. x1
= −2x1 + x2  x2
= −4x1 + 3x2 + 10 cos
t 15. x1
= 3x1 + 3x2 + 8 x2
= x1 + 5x2 + 4e3t 16. x1
= x1 + x2 + 6e3t  x2
= x1 + x2 + 4 17. x1
= 6x1 + 5x2 − 4 cos
3t x2
= x1 + 2x2 + 8 18. x1
= 3x1 − 2x2 + 3e2t  x2
= 9x1 − 3x2 + e2t In each of Problems 19 through 23, solve the initial value problem by diagonalization. 19.

x1


 x2


= x1 + x2 + 6e = x1 + x2 + 2e  x1
0 = 6 x2
0 = 0 2t

2t

20. x1
= x1 − 2x2 + 2t x2
= −x1 + 2x2 + 5 x1
0 = 13 x2
0 = 12 21. x1
= 2x1 − 5x2 + 5 sin
t x2
= x1 − 2x2  x1
0 = 10 x2
0 = 5 22. x1
= 5x1 − 4x2 + 4x3 − 3e−3t  x2
= 12x1 − 11x2 + 12x3 + t x3
= 4x1 − 4x2 + 5x3  x1
0 = 1 x2
0 = −1 x3
0 = 2 23. x1
= 3x1 − x2 − x3  x2
= x1 + x2 − x3 + t x3
= x1 − x2 + x3 + 2et  x1
0 = 1 x2
0 = 2 x3
0 = −2

CHAPTER

11

NONLINEAR SYSTEMS AND EXISTENCE OF SOLU TIONS THE PHASE PLANE, PHASE PORTRAITS, AND DIRECTION FIELDS CRITICAL POINTS AND STABILITY ALMOST LINEAR SYSTEMS PREDATOR/PREY POP

Qualitative Methods and Systems of Nonlinear Differential Equations

11.1

Nonlinear Systems and Existence of Solutions The preceding chapter was devoted to matrix methods for solving systems of differential equations. Matrices are suited to linear problems. In algebra, the equations we solve by matrix methods are linear, and in differential equations the systems we solve by matrices are also linear. However, many interesting problems in mathematics, the sciences, engineering, economics, business and other areas involve systems of nonlinear differential equations, or nonlinear systems. We will consider such systems having the special form: x1

t = F1
t x1  x2      xn  x2

t = F2
t x1  x2      xn   

(11.1)

xn

t = Fn
t x1  x2      xn  This assumes that each equation of the system can be written with a first derivative isolated on one side, and a function of t and the unknown functions x1
t     xn
t on the other. Initial conditions for this system have the form x1
t0  = x10  x2
t0  = x20      xn
t0  = xn0 

(11.2)

in which t0 is a given number and x10      xn0 are given numbers. An initial value problem consists of finding a solution of the system (11.1) satisfying the initial conditions (11.2). We will state an existence/uniqueness result for this initial value problem. In the statement, an open rectangular parallelepiped in
n + 1− dimensional t x1      xn − space consists of all points
t x1      xn  in Rn+1 whose coordinates satisfy inequalities  < t <  1 < x1 < 1      n < xn < n  403

404

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations

If n = 1 this is an open rectangle in the t x plane, and in three space these points form an open three-dimensional box in 3-space. “Open” means that only points in the interior of parallelopiped, and no points on the bounding faces, are included. THEOREM 11.1

Existence/Uniqueness for Nonlinear Systems

Let F1      Fn and their first partial derivatives be continuous at all points of an open rectangular parallelepiped K in Rn+1 . Let
t0  x10      xn0  be a point of K. Then there exists a positive number h such that the initial value problem consisting of the system (11.1) and the initial conditions (11.2) has a unique solution x1 = 1
t x2 = 2
t     xn = n
t defined for t0 − h < t < t0 + h. Many systems we encounter are nonlinear and cannot be solved in terms of elementary functions. This is why we need an existence theorem, and why we will shortly develop qualitative methods to determine properties of solutions without having them explicitly in hand. As we develop ideas and methods for analyzing nonlinear systems, it will be helpful to have some examples to fall back on and against which to measure new ideas. Here are two examples that are important and that come with some physical intuition about how solutions should behave.

EXAMPLE 11.1 The Simple Damped Pendulum

We will derive a system of differential equations describing the motion of a simple pendulum, as shown in Figure 11.1. Although we have some intuition about how a pendulum bob should move, nevertheless the system of differential equations describing this motion is nonlinear and cannot be solved in closed form. Suppose the pendulum bob has mass m, and is at the end of a rod of length L. The rod is assumed to be so light that its weight does not figure into the motion of the bob. It serves only to constrain the bob to remain at fixed distance L from the point of suspension. The position of the bob at any time is described by its displacement angle 
t from the vertical. At some time we call t = 0 the bob is displaced by an angle 0 and released from rest. To describe the motion of the bob, we must analyze the forces acting on it. Gravity acts downward with a force of magnitude mg. The damping force (air resistance, friction of the bar at its pivot point) is assumed to have magnitude c

t for some positive constant c. By Newton’s laws of motion, the rate of change of angular momentum about any point, with respect to time,

cL θ'(t)

θ

m mg j

L 1
cos(θ)


L sin(θ) i FIGURE 11.1

pendulum.

Simple, damped

11.1 Nonlinear Systems and Existence of Solutions

405

equals the moment of the resultant force about that point. The angular momentum is mL2 

t. From the diagram, the horizontal distance between the bob and the vertical center position at time t is L sin

t. Then, mL2 


t = −cL

t − mgL sin

t The negative signs on the right take into account the fact that, if the bob is displaced to the right, these forces tend to make the bob rotate clockwise, which is the negative orientation. It is customary to write this differential equation as 

+ !
+ 2 sin
 = 0

(11.3)

with ! = c/mL and  = g/L. Convert this second order equation to a system as follows. Let 2

x =  y = 
 Then the pendulum equation (11.3) becomes x
= y y
+ !y + 2 sin
x = 0 or x
= y y
= −2 sin
x − !y This is a nonlinear system because of the sin
x term. We cannot write a solution of this system in closed form. However, we will soon have methods to analyze the behavior of solutions and hence the motion itself. In matrix form, the pendulum system is     0 1 0 X+  X
= 0 −! −2 sin
x   x in which X = . y

EXAMPLE 11.2 Nonlinear Spring

Consider an object of mass m attached to a spring. If the object is displaced and released, its motion is governed by Hooke’s law, which states that the force exerted on the mass by the spring is F
r = −kr, with k a positive constant and r the distance displaced from the equilibrium position (position at which the object is at rest). Figure 2.4 shows a typical such mass/spring system, with r used here for the displacement instead of y used in Chapter 2. This is a linear model, since F is a linear function (a constant times r to the first power). The spring model becomes nonlinear if F
r is nonlinear. Simple nonlinear models are achieved by adding terms to −kr. What kind of terms should we add? Intuition tells us that the spring should not care whether we displace an object left or right before releasing it. Since displacements in opposite directions carry opposite signs, this means that we want F
−r = −F
r, so F should be an odd function.This suggests adding multiples of odd powers of r. The simplest such model is F
r = −kr + r 3 

CHAPTER 11

406

Qualitative Methods and Systems of Nonlinear Differential Equations

If we also allow a damping force which in magnitude is proportional to the velocity, then by Newton’s law this spring motion is governed by the second order differential equation mr

= −kr + r 3 − cr
 To convert this to a system, let x = r and y = r
. The system is x
= y y
= − In matrix form, this system is  X
=

SECTION 11.1

0 −k/m

1 −c/m



 X+

0 x3 /m

 

PROBLEMS x
0 = a, y
0 = b. What are the physical interpretations of the initial conditions? Are there any restrictions on the numbers a and b in applying the theorem to assert the existence of a unique solution in some interval
−h h?

1. Apply the existence/uniqueness theorem to the system for the simple damped pendulum, with initial conditions x
0 = a, y
0 = b. What are the physical interpretations of the initial conditions? Are there any restrictions on the numbers a and b in applying the theorem to assert the existence of a unique solution in some interval
−h h?

3. Suppose the driving force for the nonlinear spring has additional terms, say F
r = −kr + r 3 + r 5 . Does this problem still have a unique solution in some interval
−h h?

2. Apply the existence/uniqueness theorem to the system for the nonlinear spring system, with initial conditions

11.2

 k c x + x3 − y m m m

The Phase Plane, Phase Portraits and Direction Fields Throughout this chapter we will consider systems of two first-order differential equations in two unknowns. In this case it is convenient to denote the variables as x and y rather than x1 and x2 . Thus consider the system x

t = f
x
t y
t y

t = g
x
t y
t

(11.4)

in which f and g are continuous, with continuous first partial derivatives, in some part of the plane. We often write this system as X
= F
x
t y
t where

 X=

x y



 and F
x y =

f
x y g
x y

 

The system (11.4) is a special case of the system (11.1). We assume in (11.4) that neither f nor g has an explicit dependence on t. Rather, f and g depend only on x and y, and t appears only through dependencies of these two variables on t. We refer to such a system as autonomous.

11.2 The Phase Plane, Phase Portraits and Direction Fields

407

Working in the plane will allow us the considerable advantage of geometric intuition. If x = 
t, y = 
t is a solution of (11.4), the point

t 
t traces out a curve in the plane as t varies. Such a curve is called a trajectory, or orbit, of the system. A copy of the plane containing drawings of trajectories is called a phase portrait for the system (11.4). In this context, the x y plane is called the phase plane. We may consider trajectories as oriented, with

t 
t moving along the trajectory in a certain direction as t increases. If we think of t as time, then

t 
t traces out the path of motion of a particle, moving under the influence of the system (11.4), as time increases. In the case of orbits that are closed curves, we take counterclockwise orientation as the positive orientation, unless specific exception is made. In some phase portraits, short arrows are also drawn. The arrow at any point is along the tangent to the trajectory through that point, and in the direction of motion along this trajectory. This type of drawing combines the phase portrait with a direction field, and gives an overall sense of the flow of the trajectories, as well as graphs of some specific trajectories. One way to construct trajectories is to write dy dy/dt g
x y = =  dx dx/dt f
x y Because the system is autonomous, this is a differential equation in x and y and we can attempt to solve it and graph solutions. If the system is nonautonomous, then f/g may depend explicitly on t then we cannot use this strategy to generate trajectories.

EXAMPLE 11.3

Consider the autonomous system x
= y = f
x y y
= x2 y2 = g
x y Then x2 y 2 dy = = x2 y dx y a separable differential equation we write as 1 dy = x2 dx y Integrate to get 1 ln y = x3 + C 3 or 3

y = Aex /3  Graphs of these curves for various values of A form trajectories of this system, some of which are shown in Figure 11.2.

408

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y 6 4 2


3


2


1

0

2

1

x


2
4
6  FIGURE 11.2

Some trajectories of the system

x
= y y
= x2 y 2

.

EXAMPLE 11.4

For the autonomous system x
= −2y − x sin
xy y
= 2x + y sin
xy

(11.5)

we have dy 2x + y sin
xy =−  dx 2y + x sin
xy This is not separable, but we can write
2x + y sin
xydx +
2y + x sin
xydy = 0 which is exact. We find the potential function H
x y = x2 + y2 − cos
xy, and the general solution of this differential equation is defined implicitly by H
x y = x2 + y2 − cos
xy = C in which C is an arbitrary constant. Figure 11.3 shows a phase portrait for this system (11.5), consisting of graphs of these curves for various choices of C. Usually we will not be so fortunate as to be able to solve dy/dx = g
x y/f
x y in closed form. In such a case we may still be able to use a software package to generate a phase portrait. Figure 11.4 is a phase portrait of the system x
= x cos
y y
= x2 − y3 + sin
x − y

11.2 The Phase Plane, Phase Portraits and Direction Fields

409

y 2

1

–2

–1

1

0

2

x

–1

–2  FIGURE 11.3 Phase portrait for

x
= −2y − x sin
xy y
= 2x + y sin
xy

.

y 2

1

0

0.6

0.8

1.0

1.2

1.4

1.6

1.8

2.0

2.2

x

–1

–2

–3  FIGURE 11.4 Phase portrait for

x
= x cos
y y
= x2 − y3 + sin
x − y

.

generated in this way. Figure 11.5 (p. 410) is a phase portrait for a damped pendulum with 2 = 10 and ! = 03, and Figure 11.6 (p. 410) is a phase portrait for a nonlinear spring system with  = 02, k/m = 4 and c/m = 2. We will consider phase portraits for the damped pendulum and nonlinear spring system in more detail when we treat almost linear systems. If x = 
t, y = 
t is a solution of (11.4), and c is a constant, we call the pair 
t + c, 
t + c a translation of  and . We will use the following fact.

410

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y

8 6 4 2
15


10


5

5

10

15

x


2
4
6 FIGURE 11.5

Phase portrait for a damped pendulum.

y 30 20 10
10


5

0

5

10

x


10
20
30
40 FIGURE 11.6

Phase portrait for a nonlinear spring.

LEMMA 11.1

A translation of a solution of the system (11.4) is also a solution of this system. Proof

Suppose x = 
t, y = 
t is a solution. This means that x

t = 

t = f

t 
t and y

t = 

t = g

t 
t

Let , x
t = 
t + c and , y
t = 
t + c

11.2 The Phase Plane, Phase Portraits and Direction Fields

411

for some constant c. By the chain rule, d, x d
t + c d
t + c = = 

t + c dt d
t + c dt = f

t + c 
t + c = f
, x
t, y
t and, similarly, d, y = 

t + c = g

t + c 
t + c = g
, x
t, y
t dt Therefore x = , x
t, y =, y
t is also a solution. We may think of a translation as a reparamatrization of the trajectory, which of course does not alter the fact that it is a trajectory. If we think of the point

t 
t as moving along the orbit, a translation simply means rescheduling the point to change the times at which it passes through given points of the orbit. We will need the following facts about trajectories.

THEOREM 11.2

Let F and G be continuous, with continuous first partial derivatives, in the
x y plane. Then, 1. If
a b is a point in the plane, there is a trajectory through
a b 2. Two trajectories passing through the same point must be translations of each other. Proof

Conclusion (1) follows immediately from Theorem 11.1, since the initial value problem x
= f
x y y
= g
x y

x
0 = a y
0 = b

has a solution, and the graph of this solution is a trajectory through
a b. For (2), suppose x = 1
t, y = 1
t and x = 2
t, y = 2
t are trajectories of the system (11.4). Suppose both trajectories pass through
a b. Then for some t0 , 1
t0  = a and 1
t0  = b and for some t1 , 2
t1  = a and 2
t1  = b Let c = t0 − t1 and define , x
t = 1
t + c and , y
t = 1
t + c. Then x = , x
t, y = , y
t is a trajectory, by Lemma 11.1. Further, , x
t1  = 1
t0  = a and , y
t1  = 
t0  = b Therefore x = , x
t, y =, y
t is the unique solution of the initial value problem x
= f
x y y
= g
x y

x
t1  = a y
t1  = b

But x = 2
t, y = 2
t is also the solution of this problem. Therefore, for all t, 2
t = , x
t = 1
t + c and 2
t =, y
t = 1
t + c This proves that the two trajectories x = 1
t, y = 1
t and x = 2
t, y = 2
t are translations of each other.

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If we think of translations of trajectories as the same trajectory (just a change in the parameter), then conclusion (2) states that distinct trajectories cannot cross each other. This would violate uniqueness of the solution of the system that passes through the point of intersection. Conclusion (2) of Theorem 11.2 does not hold for systems that are not autonomous.

EXAMPLE 11.5

Consider the system 1 x

t = x = f
t x y t 1 y

t = − y + x = g
t x y t This is nonautonomous, since f and g have explicit t-dependencies. We can solve this system. The first equation is separable. Write 1 1 dx = dt x t to obtain x
t = ct. Substitute this into the second equation to get 1 y
+ y = ct t a linear first-order differential eqaution. This equation can be written ty
+ y = ct2  or
ty
= ct2  Integrate to get c ty = t3 + d 3 Hence c 1 y
t = t2 + d  3 t Now observe that conclusion (2) of Theorem 11.2 fails for this system. For example, for any number t0 , the trajectory x
t =

1 1 2 t02 1  t y
t = t − t0 3t0 3 t

passes through
1 0 at time t0 . Because t0 is arbitrary, this gives many trajectories passing through
1 0 at different times, and these trajectories are not translations of each other. We now have the some of the vocabulary and tools needed to analyze 2 × 2 nonlinear autonomous systems of differential equations. First, however, we will reexamine linear systems, which we know how to solve explicitly. This will serve two purposes. It will give us some experience with phase portraits, as well as insight into significant features that solutions of a system might have. In addition, we will see shortly that some nonlinear systems can be thought of as perturbations of linear systems (that is, as linear systems with “small” nonlinear terms added). In such a case, knowledge of solutions of the linear system yields important information about solutions of the nonlinear system.

11.3 Phase Portraits of Linear Systems

SECTION 11.2

413

PROBLEMS

In each of Problems 1 through 6, find the general solution of the system and draw a phase portrait containing at least six trajectories of the system. 1. x
= 4x + y y
= −17x − 4y

7. x
= 9y y
= −4x 8. x
= 2xy y
= y2 − x2 9. x
= y + 2 y
= x − 1

2. x
= 2x y
= 8x + 2y

10. x
= csc
x y
= y

3. x
= 4x − 7y y
= 2x − 5y

11. x
= x y
= x + y

4. x
= 3x − 2y y
= 10x − 5y

12. x
= x2  y
= y

5. x
= 5x − 2y y
= 4y 13. How would phase portraits for the following systems compare with each other?

6. x
= −4x − 6y y
= 2x − 11y In each of Problems 7 through 12, use the method of Examples 11.3 and 11.4 to draw some integral curves (at least six) for the system.

11.3

(a) x
= F
x y y
= G
x y (b) x
= −F
x y y
= −G
x y

Phase Portraits of Linear Systems In preparation for studying the nonlinear autonomous system (11.4), we will thoroughly analyze the linear system (11.6) X
= AX   x in which A is a 2 × 2 real matrix and X = . We assume that A is nonsingular, so the y equation AX = O has only the trivial solution. For the linear system X
= AX, we actually have the solutions in hand. We will examine these solutions to prepare for the analysis of nonlinear systems, for which we are unlikely to have explicit solutions. The origin
0 0 stands apart from other points in the plane in the following respect. The trajectory through the origin is the solution of:   0
X = AX X
0 = O = 0 and this is the constant trajectory x
t = 0 y
t = 0 for all t. The graph of this trajectory is the single point
0 0. For this reason, the origin is called an equilibrium point of the system, and the constant solution X = O is called an equilibrium solution. The origin is also called a critical point of X
= AX. By Theorem 11.1, no other trajectory can pass through this point. As we proceed, observe how the behavior of trajectories of X
= AX near this critical point is the key to understanding the behavior of trajectories throughout the entire plane. The critical point, then, will be the focal point in drawing a phase portrait of the system and analyzing the behavior of solutions.

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Qualitative Methods and Systems of Nonlinear Differential Equations

We will draw the phase portrait for X
= AX in all cases that can occur. Because the general solution of (11.6) is completely determined by the eigenvalues of A, we will use these eigenvalues to distinguish cases. Case 1—Real, distinct eigenvalues  and  of the same sign. Let associated eigenvectors be, respectively, E1 and E2 . Since  and  are distinct, E1 and E2 are linearly independent. The general solution is   x
t X
t = = c1 E1 et + c2 E2 et  y
t Since E1 and E2 are vectors in the plane, we can represent them as arrows from the origin, as in Figure 11.7. Draw half-lines L1 and L2 from the origin along these eigenvectors, respectively, as shown. These half-lines lines are parts of trajectories, and so do not pass through the origin, which is itself a trajectory. y L2 L1

E2 E1

x

FIGURE 11.7 Eigenvectors E1 and E2 of X
= AX, for distinct eigenvalues of A.

Now consider subcases. Case 1-(a)—The eigenvalues are negative, say  <  < 0. Since et → 0 and et → 0 as t → , then X
t →
0 0 and every trajectory approaches the origin as t → . However, this can happen in three ways, depending on an initial point P0 
x0  y0  we choose for a trajectory to pass through at time t = 0. Here are the three possibilities. If P0 is on L1 , then c2 = 0 and X
t = c1 E1 et  which for any t is a scalar multiple of E1 . The trajectory through P0 is the half-line from the origin along L1 through P0 , and the arrows toward the origin indicate that points on this trajectory approach the origin along L1 as t increases. This is the trajectory T1 of Figure 11.8. If P0 is on L2 , then c1 = 0 and now X
t = c2 E2 et . This trajectory is a half-line from the origin along L2 through P0 . Again, the arrows indicate that points on this trajectory also approach the origin along L2 as t → . This is the trajectory T2 of Figure 11.8. If P0 is on neither L1 or L2 , then the trajectory is a curve through P0 having the parametric form X
t = c1 E1 et + c2 E2 et  Write this as

 X
t = et c1 E1 e
−t + c2 E2 

11.3 Phase Portraits of Linear Systems

415

y T3

P0

L2

L1

P0 T1

T2

P0 x

FIGURE 11.8 Trajectories along E1 , along E2 , or asymptotic to E2 in the case  <  < 0.

Because  −  < 0, e
−t → 0 as t → and the term c1 E1 e
−t exerts increasingly less influence on X
t. In this case, X
t still approaches the origin, but also approaches the line L2 , as t → . A typical such trajectory is shown as the curve T3 of Figure 11.8. A phase portrait of X
= AX in this case therefore has some trajectories approaching the origin along the lines through the eigenvectors of A and all others approaching the origin along curves that approach one of these lines asymptotically. In this case the origin is called a nodal sink of the system X
= AX. We can think of particles flowing along the trajectories and toward the origin. The following example and phase portrait are typical of nodal sinks.

EXAMPLE 11.6

Consider the system X
= AX, in which A=



−6 5

−2 1

 

A has engenvalues and corresponding eigenvectors     2 −1 −1 and − 4  −5 1 In the notation of the discussion,  = −4 and  = −1. The general solution is     −1 2 −4t e + c2 e−t  X
t = c1 1 −5 L1 is the line through
−1 1, and L2 the line through
2 −5. Figure 11.9 shows a phase portrait for this system, with the origin a nodal sink. Case 1-(b)—The eigenvalues are positive, say 0 <  < . The discussion of Case 1-(a) can be replicated with one change. Now et and et approach instead of zero as t increases. The phase portrait is like that of the previous case, except all the arrows are reversed and trajectories flow away from the origin instead of into the origin as time increases. As we might expect, now the origin is called a nodal source. Particles are flowing away from the origin. Here is a typical example of a nodal source.

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2


3


2


1

0

1

2

x


2


4 Phase portrait showing a nodal sink x
= −6x − 2y . y
= 5x + y

FIGURE 11.9 

of

EXAMPLE 11.7

Consider the system

 X
=

3 1

3 5

 

This has eigenvalues and eigenvectors     −3 1 2 and 6  1 1 Now  = 6 and  = 2, and the general solution is     −3 1 2t e + c2 e6t  X
t = c1 1 1 Figure 11.10 shows a phase portrait for this system, exhibiting the behavior expected for a nodal source at the origin. Case 2—Real, distinct eigenvalues of opposite sign. Suppose the eigenvalues are  and  with  < 0 < . The general solution still has the appearance X
t = c1 E1 et + c2 E2 et  and we start to draw a phase portrait by again drawing half-lines L1 and L2 from the origin along the eigenvectors. If P0 is on L1 , then c2 = 0 and X
t moves on this half-line away from the origin as t increases, because  > 0 and et → as t → . But if P0 is on L2 , then c1 = 0 and X
t moves along this half-line toward the origin, because  < 0 and et → 0 as t → .

11.3 Phase Portraits of Linear Systems

417

y 4

2


3


2


1

0

1

2

x


2


4 

x
= 3x + 3y

FIGURE 11.10 Nodal source of

y
= x + 5y

.

The arrows along the half-lines along the eigenvectors therefore have opposite directions, toward the origin along L2 and away from the origin along L1 . This is in contrast to Case 1, in which solutions starting out on the half-lines through the eigenvectors either both approached the origin or both moved away from the origin as time increased. If P0 is on neither L1 nor L2 , then the trajectory through P0 does not come arbitrarily close to the origin for any times, but rather approaches the direction determined by the eigenvector E1 as t → (in which case et → 0) or the direction determined by E2 as t → − (in which case et → 0). The phase portrait therefore has typical trajectories as shown in Figure 11.11. The

y 4

2


4


2

0

2

4


2


4

FIGURE 11.11 Typical phase portrait for a saddle point at the origin.

x

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lines along the eigenvectors determine four trajectories that separate the plane into four regions. A trajectory starting in one of these regions must remain in it because distinct trajectories cannot cross each other, and such a trajectory is asymptotic to both of the lines bounding its region. The origin in this case is called a saddle point.

EXAMPLE 11.8

Consider X
= AX with



−1 3 2 −2

A=

 

Eigenvalues and eigenvectors of A are     −1 3 −4 and 1 1 2 The general solution is  X
t = c1

−1 1



 e−4t + c2



3 2

et

and a phase portrait is given in Figure 11.12. In this case of a saddle point at the origin, trajectories do not enter or leave the origin, but asymptotically approach the lines determined by the eigenvectors. Case 3—Equal eigenvalues. Suppose A has the real eigenvalue  of multiplicity 2. There are two possibilities.

y

150

100

50


150


100


50

0

50

 FIGURE 11.12

Saddle point of

x
= −x + 3y y
= 2x − 2y

.

x

11.3 Phase Portraits of Linear Systems

419

Case 3-(a)—A has two linearly independent eigenvectors E1 and E2 . Now the general solution of X
= AX is X
t =
c1 E1 + c2 E2  et  If

 E1 =

a b



 and E2 =

h k



then, in terms of components, x
t =
c1 a + c2 het 

y
t =
c1 b + c2 ket 

Now y
t = constant. x
t This means that all trajectories in this case are half-lines from the origin. If  > 0, arrows along these trajectories are away from the origin, as in Figure 11.13. If  < 0, they move toward the origin, reversing the arrows in Figure 11.13. y

x

FIGURE 11.13 Typical proper node with positive eigenvalue of A.

The origin in case 3 - (a) is called a proper node. Case 3-(b)—A does not have two linearly independent eigenvectors. In this case there is an eigenvector E and the general solution has the form X
t = c1
Et + Wet + c2 Eet =
c1 W + c2 E + c1 Et et  To visualize the trajectories, begin with arrows from the origin representing the vectors W and E. Now, for selected constants c1 and c2 , draw the vector c1 W + c2 E, which may have various orientations relative to W and E, depending on the signs and magnitudes of c1 and c2 . Some possibilities are displayed in Figure 11.14. For given c1 and c2 , the vector c1 W + c2 E + c1 Et drawn as an arrow from the origin, sweeps out a straight line L as t varies over all real values. For a given t, X
t is the vector c1 W + c2 E + c1 Et from the origin to a point on L, with length adjusted by a factor et . If  is negative, then this length goes to zero as t → and the vector X
t sweeps out a curve as shown in Figure 11.15, approaching the origin tangent to E. If

420

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W  2E

W W
E E

x


W  2E FIGURE 11.14 Vectors c1 W + c2 E in the case of an improper node.

y L c1W  c2E  c1Et X(t) = [c1W  c2E  c1Et]eλt W

FIGURE 11.15

E

x

Typical trajectory near an improper

node.

 > 0, we have the same curve (now et → 0 as t → − ), except that the arrow indicating direction of the flow on the trajectory is reversed. The origin in this case is called a improper node of the system X
= AX. The following example has a phase portrait that is typical of improper nodes.

EXAMPLE 11.9

Let  A=

−10 −6



6 2

 

Then A has eigenvalue −4, and every eigenvector is a real constant multiple of E = A routine calculation gives  W=

1 7 6

 

1 1

 .

11.3 Phase Portraits of Linear Systems

421

and the general solution is  X
t = c1

t+1 t + 76



 e−4t + c2

1 1

 e−4t 

Figure 11.16 is a phase portrait for this system. We can see that the trajectories approach the origin tangent to E in this case of an improper node at the origin, with negative eigenvalue for A.

y 60 50 40 30 20 10
100


50

0

50

100

x


10

FIGURE 11.16 Phase portrait for the improper node



of

x
= −10x + 6y y
= −6x + 2y

.

Case 4—Complex eigenvalues with nonzero real part. We know that the complex eigenvalues must be complex conjugates, say  =  + i and  =  − i. The complex eigenvectors are also conjugates. Write these, respectively, as U + iV and U − iV. Then the general solution of X
= AX is X
t = c1 et U cos
t − V sin
t + c2 et U sin
t + V cos
t Suppose first that  < 0. The trigonometric terms in this solution cause X
t to rotate about the origin as t increases, while the factor et causes X
t to move closer to the origin (or, equivalently, the length of the vector X
t to decrease to zero) as t → . This suggests a trajectory that spirals inward toward the origin as t increases. Since t varies over the entire real line, taking on both negative and positive values, the trajectories when  > 0 have the same spiral appearance, but now the arrows are reversed and X
t moves outward, away from the origin, as t → . The origin in this case is called a spiral point. When  < 0 the origin is a spiral sink because the flow defined by the trajectories is spiralling into the origin. When  > 0 the origin is a spiral source because now the origin appears to be spewing material outward in a spiral pattern. The phase portrait in the following example is typical of a spiral source.

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EXAMPLE 11.10

Let

 A=

−1 4

−2 3

  

with eigenvalues 1 + 2i and 1 − 2i and eigenvectors, respectively, Let

 U=

−1 2



 and V =

1 0

−1 + i 2



 and

−1 − i 2

 .



so that the eigenvectors are U + iV and U − iV. The general solution of X
= AX is    

−1 1 t X
t = c1 e cos
2t − sin
2t 2 0    

−1 1 + c2 et sin
2t + cos
2t  2 0 Figure 11.17 gives a phase portrait for this system, showing trajectories spiraling away from the spiral source at the origin because the real part of the eigenvalues is positive.

y

40

20


40


30


20


10

0

10

20

x


20

FIGURE 11.17

Spiral source of the system

 x
= −x − 2y y
= 4x + 3y

.

Case 5—Pure imaginary eigenvalues. Now trajectories have the form X
t = c1 U cos
t − V sin
t + c2 U sin
t + V cos
t Because of the trigonometric terms, this trajectory moves about the origin. Unlike the preceding case, however, there is no exponential factor to decrease or increase distance from the origin as t increases. This trajectory is a closed curve about the origin, representing a periodic solution

11.3 Phase Portraits of Linear Systems

423

of the system. The origin in this case is called a center of X
= AX. In general, any closed trajectory of X
= AX represents a periodic solution of this system.

EXAMPLE 11.11

Let



3 18 −1 −3

A=

  

−3 − 3i 1



A has eigenvalues 3i and −3i, with respective eigenvectors, eigenvectors and   −3 + 3i . A phase portrait is given in Figure 11.18, showing closed trajectories about the 1 center (origin). If we wish, we can write the general solution y

10

5


40


20

0

20

40

x


5


10

FIGURE 11.18 Center of the system

 x
= 3x + 18y y
= −x − 3y

.



  

−3 3 cos
3t + sin
3t 1 0    

−3 −3 + c2 sin
3t + cos
3t  1 0

X
t = c1

We now have a complete description of the behavior of trajectories for the 2 × 2 constant coefficient system X
= AX. The general appearance of the phase portrait is completely determined by the eigenvalues of A, and the critical point
0 0 is the primary point of interest, with the following correspondences: Real, distinct eigenvalues of the same sign—
0 0 is a nodal source (Figure 11.10, p. 417) or sink (Figure 11.9, p. 416). Real, distinct eigenvalues of opposite sign—
0 0 is a saddle point (Figure 11.12, p. 418).

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424

Qualitative Methods and Systems of Nonlinear Differential Equations

Equal eigenvalues, two linearly independent eigenvectors—
0 0 is a proper node (Figure 11.13, p. 419). Equal eigenvalues, all eigenvectors a multiple of a single eigenvector—
0 0 is an improper node (Figure 11.16, p. 421) Complex eigenvalues with nonzero real part—
0 0 is a spiral point (Figure 11.17, p. 422). Pure imaginary eigenvalues—
0 0 is a center (Figure 11.18, p. 423). When we speak of a classification of the origin of a linear system, we mean a determination of the origin as a nodal source or sink, saddle point, proper or improper node, spiral point or center.

SECTION 11.3

PROBLEMS 4. x
= 9x − 7y y
= 6x − 4y

In each of Problems 1 through 10, use the eigenvalues of the matrix of the system to classify the origin of the system. Draw a phase portrait for the system. It is assumed here that software is available to do this, and it is not necessary to solve the system to generate the phase portrait.

5. x
= 7x − 17y y
= 2x + y 6. x
= 2x − 7y y
= 5x − 10y 7. x
= 4x − y y
= x + 2y

1. x
= 3x − 5y y
= 5x − 7y

8. x
= 3x − 5y y
= 8x − 3y

2. x
= x + 4y y
= 3x

9. x
= −2x − y y
= 3x − 2y

3. x
= x − 5y y
= x − y

11.4

10. x
= −6x − 7y y
= 7x − 20y

Critical Points and Stability A complete knowledge of the possible phase portraits of linear 2 ×2 systems is good preparation for the analysis of nonlinear systems. In this section we will introduce the concept of critical point for a nonlinear system, define stability of critical points, and prepare for the qualitiative analysis of nonlinear systems, in which we attempt to draw conclusions about how solutions will behave, without having explicit solutions in hand. We will consider the 2 × 2 autonomous system x

t = f
x
t y
t y

t = g
x
t y
t or, more compactly, x
= f
x y y
= g
x y This is the system (11.4) discussed in Section 11.2. We will assume that f and g are continuous with continuous first partial derivatives in some region D of the x y− plane. In specific cases D may be the entire plane. This system can be written in matrix form as X
= F
X

11.4 Critical Points and Stability

425

In August of 1999, the Petronas Towers was officially opened. Designed by the American firm of Cesar Pelli and Associates, in collaboration with Kuala Lumpur City Center architects, the graceful towers have an elegant slenderness (height to width) ratio of 9:4. This was made possible by modern materials and building techniques, featuring high-strength concrete that is twice as effective as steel in sway reduction. The towers are supported by 75-foot-by-75-foot concrete cores and an outer ring of super columns. The 88 floors stand 452 meters above street level, and include 65,000 square meters of stainless steel cladding and 77,000 square meters of vision glass. Computations of stability of structures involve the analysis of critical points of systems of nonlinear differential equations.

where

 X
t =

x
t y
t



 and F
X =

f
x y g
x y

 

Taking the lead from the linear system X
= AX, we make the following definition.

DEFINITION 11.1

Critical Point

A point
x0  y0  in D is a critical point (or equilibrium point) of X
= F
X if f
x0  y0  = g
x0  y0  = 0

We see immediately one significant difference between the linear and nonlinear cases. The linear system X
= AX, with A nonsingular, has exactly one critical point, the origin. A nonlinear system X
= F
X can have any number of critical points. We will, however, only consider systems in which critical points are isolated. This means that, about any critical point, there is a circle that contains no other critical point of the system.

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EXAMPLE 11.12

Consider the damped pendulum (Example 11.1), whose motion is governed by the system x
= y y
= −2 sin
x − !y Here f
x y = y and g
x y = −2 sin
x − !y The critical points are solutions of y=0 and −2 sin
x − !y = 0 These equations are satisfied by all points
n 0, in which n = 0 ±1, ±2,    . These critical points are isolated. About any point
n 0, we can draw a circle (for example, of radius 1/4) that does not contain any other critical point. For this problem, the critical points split naturally into two classes. Recall that x =  is the angle of displacement of the pendulum from the vertical downward position, with the bob at the bottom, and y = d/dt. When n is even, then x =  = 2k for k any integer. Each critical point
2k 0 corresponds to the bob pointing straight down, with zero velocity (because y = x
= 
= 0). When n is odd, then x =  =
2k + 1 for k any integer. The critical point

2k + 1 0 corresponds to the bob in the vertical upright position, with zero velocity. Without any mathematical analysis, there is an obvious and striking difference between these two kinds of critical points. At, for example, x = 0 , the bob hangs straight down from the point of suspension. If we displace it slightly from this position and then release it, the bob will go through some oscillations of decreasing amplitude, after which it will return to its downward position and remain there. This critical point, and all critical points
2n 0, are what we will call stable. Solutions of the pendulum equation for initial values near this critical point remain close to the constant equilibrium solution for all later times. By contrast, consider the critical point
 0. This has the bob initially balanced vertically upward. If the bob is displaced, no matter how slightly, it will swing downward and oscillate back and forth some number of times, but never return to this vertical position. Solutions near this constant equilibrium solution (bob vertically up) do not remain near this position, but move away from it. This critical point, and any critical point

2k + 1 0, is unstable.

EXAMPLE 11.13

Consider the damped nonlinear spring of Example 11.2. The system of differential equations governing the motion is x
= y y
= −

k c  x + x3 − y m m m

11.4 Critical Points and Stability

427

√ √ The critical points are
0 0,
k/ 0 and
− k/ 0. Recall that x measures the position of the spring, from the equilibrium (rest) position, and y = dx/dt is the velocity of the spring. We will do a mathematical analysis of this system shortly, but for now look at these critical points from the point of view of our experience with how springs behave. If we displace the spring very slightly from the equilibrium solution
0 0, and then release it, we expect it to undergo some motion back and forth and then come to rest, approaching the equilibrium point. In this sense
0 0 is a stable critical √ point. However, if we displace the spring slightly from a position very nearly at distance k/ to the right or left of the equilibrium position and then release it, the spring may or may not return to this position, depending on the relative sizes of the damping constant c and the coefficients in the nonlinear spring force function, particularly . In this sense these equilibrium points may be stable or may not be. In the next section we will develop the tools for a more definitive analysis of these critical points. Taking a cue from these examples, we will define a concept of stability of critical points. Recall that  V = v1 2 + v2 2 is the length (or norm) of a vector V =
v1  v2  in the plane. If W =
w1  w2  is also a vector in the plane, then V − W is the length of the vector from W to V. If W =
w1  w2 , then  1/2 V − W =
v1 − w1 2 +
v2 − w2 2 is also the distance between the points
v1  v2  and
w1  w2 . Finally, if X0 is a given vector, then the locus of points (vectors) X such that X − X0  < r for any positive r, is the set of points X within the circle of radius r about X0 . These are exactly the points at distance 0. 2 −1 +  (a) Show that, when  = 0, the critical point is a center, stable but not asymptotically stable. Generate a phase portrait for this system. (b) Show that, when  = 0, the critical point is not a center, no matter how small  is chosen. Generate a 1 phase portrait for this system with  = . 10 This problem illustrates the sensitivity of trajectories of the system to small changes in the coefficients, in the case of pure imaginary eigenvalues.

11.5

12. Consider the system X
= AX, where A =  2+ 5 and  > 0. −5 −8 (a) Show that, when  = 0, A has equal eigenvalues and does not have two linearly independent eigenvectors. Classify the type of critical point at the origin and its stability characteristics. Generate a phase portrait for this system. (b) Show that, if  is not zero (but can be arbitrarily small in magnitude), then A has real and distinct eigenvalues. Classify the type of critical point at the origin in this case, as well as its stability characteristics. 1 Generate a phase portrait for the case  =  10 This problem illustrates the sensitivity of trajectories to small changes in the coefficients, in the case of equal eigenvalues.

Almost Linear Systems Suppose X
= F
X is a nonlinear system. We want to define a sense in which this system may be thought of as “almost linear.” Suppose the system has the special form X
= AX + G
X





(11.7)

p
x y added. Any nonlinq
x y earity of the system (11.7) is in G
X. We refer to the system X
= AX as the linear part of the system (11.7).

This is a linear system X
= AX, with another term, G
X =

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Assume that p
0 0 = q
0 0 = 0 so the system (11.7) has a critical point at the origin. The idea we want to pursue is that if the nonlinear term is “small enough”, then the behavior of solutions of the linear system X
= AX near the origin may give us information about the behavior of solutions of the original, nonlinear system near this critical point. The question is: how small is “small enough?” We will assume in this discussion that A is a nonsingular, 2 × 2 matrix of real numbers, and that p and q are continuous at least within some disk about the origin. In the following definition, we refer to partial derivatives of G, by which we mean    

G

G px py Gx = = = and Gy =  qx qy

x

y

DEFINITION 11.4

Almost Linear

The system (11.7) is almost linear in a neighborhood of
0 0 if G and its first partial derivatives are continuous within some circle about the origin, and lim

X→O

G
X = 0 X

(11.8)

This condition (11.8) means that, as X is chosen closer to the origin, G
X must become small in magnitude faster than X does. This gives a precise measure of “how small” the nonlinear term must be near the origin for the system (11.7) to qualify as almost linear. If we write       x a b p
x y X= , A= and G
X = G
x y =  y c d q
x y then the system (11.7) is x
= ax + by + p
x y y
= cx + dy + q
x y Condition (11.8) now becomes lim


xy→
00

p
x y q
x y = lim  = 0  2 2
xy→
00 x +y x2 + y 2

These limits, in terms of the components of G
X, are sometimes easier to deal with than the limit of G
X / X as X approaches the origin, although the two formulations are equivalent.

EXAMPLE 11.14

The system




X =

4 −2 1 6



 X+

−4xy −8x2 y



11.5 Almost Linear Systems

433

is almost linear. To verify this, compute −4xy and 
xy→
00 x2 + y 2 lim

−8x2 y  
xy→
00 x2 + y 2 lim

There are various ways of showing that these limits are zero, but here is a device worth remembering. Express
x y in polar coordinates by putting x = r cos
 and y = r sin
. Then −4xy 4r 2 cos
 sin
 = −4r cos
 sin
 → 0 =−  r x2 + y 2 as r → 0, which must occur if
x y →
0 0. Similarly, −8x2 y r 3 cos2
 sin
 = −8r 2 cos2
 sin
 → 0 = −8  r x2 + y 2 as r → 0.

Figure 11.22(a) shows a phase portrait of this system. For comparison, a phase portrait of the linear part X
= AX is given in Figure 11.22(b). Notice a qualitative similarity between the phase portraits near the origin. This is the rationale for the definition of almost linear systems. We will now display a correspondence between the type of critical point, and its stability properties, for the almost linear system X
= AX + G and its linear part X
= AX. The behavior is not always the same. Nevertheless, in some cases which we will identify, properties of the critical point for the linear system carry over to either the same properties for the almost linear system or, if not the same, at least to important information about the nonlinear system.

y 6

4

2

0

1

2

3

x

4


2  FIGURE 11.22(a) Phase portrait for

x
= 4x − 2y − 4xy y
= x + 6y − 8x2 y

.

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Qualitative Methods and Systems of Nonlinear Differential Equations y

20 10


6


4


2

0

2

4

6

8

x


10
20
30

FIGURE 11.22(b)

Phase portrait for the linear part of the system

of Figure 11.22(a). THEOREM 11.3

Let  and  be the eigenvalues of A. Assume that X
= AX + G is almost linear. Then the following conclusions hold for the system X
= AX + G. 1. If  and  are unequal and negative, then the origin is an asymptotically stable nodal sink of X
= AX + G. If these eigenvalues are unequal and positive, then the origin is an unstable nodal source of X
= AX + G. 2. If  and  are of opposite sign, then the origin is an unstable saddle point of X
= AX + G. 3. If  and  are complex with negative real part, then the origin is an asymptotically stable spiral point of X
= AX + G. If these eigenvalues have positive real part, then the origin is an unstable spiral point. 4. If  and  are equal and negative, then the linear system has an asymptotically stable proper or improper node, while the almost linear system has an asymptotically stable node or spiral point. If  and  are equal and positive, then the linear system has an unstable proper or improper node, while the almost linear system has an unstable node or spiral point. 5. If  and  are pure imaginary (conjugates of each other), then the origin is a center of X
= AX, but may be a center or spiral point of the almost linear system X
= AX + G. Further, in the case of a spiral point of the almost linear system, the critical point may be unstable or asymptotically stable. The only case in which the linear system fails to provide definitive information of some kind about the almost linear system is that the eigenvalues of A are pure imaginary. In this event, the linear system has a stable center, while the almost linear system can have a stable center or a spiral point which may be stable or unstable. In light of this theorem, when we ask for an analysis of a critical point of an almost linear system, we mean a determination of whether the point is an asymptotically stable nodal sink, an unstable nodal source, an unstable saddle point, an asymptotically stable spiral point or unstable spiral point, or, from (5) of the theorem, either a center or spiral point. A proof of this theorem requires some delicate analysis that we will avoid. The rest of this section is devoted to examples and phase portraits.

11.5 Almost Linear Systems

435

EXAMPLE 11.15

The system




X =

−1 −1 −1 −3



 X+

x2 y 2 x3 − y 2



√ is almost √ linear and has only one critical point,
0 0. The eigenvalues of A are −2 + 2 and −2 − 2, which are distinct and negative. The origin is an asymptotically stable nodal sink of the linear system X
= AX, and hence is also a stable and asymptotically stable nodal sink of the almost linear system. Figure 11.23 (a) and (b) shows a phase portrait of the almost linear system and its linear part, respectively. y

6

4

2


0.6


0.4


0.2

0

0.2

0.4

x


2  FIGURE 11.23(a) Phase portrait for

x
= −x − y + x2 y2 . y
= −x − 3y + x3 − y2

y

20

10


4


2

0

2

4


10
20

FIGURE 11.23(b) Phase portrait for the linear part of the system of Figure 11.23(a).

x

436

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations

EXAMPLE 11.16

The system 


X =

3 −4 6 2



 X+

x2 cos
y y3



√ is almost linear. The only critical point is
0 0. The eigenvalues of A are 25 + 21 i 95 and √ 5 − 21 i 95. The linear part has an unstable spiral point at the origin. The origin is therefore an 2 unstable spiral point of the almost linear system. Phase portraits for the given nonlinear system and its linear part are shown in Figure 11.24 (a) and (b), respectively. y

4

2


1.0


0.5

0.5

0

1.0

x


2
4

 FIGURE 11.24(a)

Phase portrait for

x
= 3x − 4y + x2 cos
y . y
= 6x + 2y + y3

y

4 2


2


1

0

1

2


2
4

FIGURE 11.24(b) Phase portrait for the linear part of the system of Figure 11.24(a).

x

11.5 Almost Linear Systems

437

EXAMPLE 11.17

The system




X =

−1 2

2 3



 X+

x sin
y 8 sin
x



√ is almost √ linear, and its only critical point is the origin. The eigenvalues of A are 1 + 2 2 and 1 − 2 2, which are real and of opposite sign. The origin is an unstable saddle point of the linear part, hence also of the given system. Phase portraits of both systems are shown in Figure 11.25 (a) (nonlinear system) and (b) (linear part).

EXAMPLE 11.18

The system



4 −2




X =



11 −4

 X+

x sin
y sin
y



√ √ is almost linear, and its only critical point is
0 0. The eigenvalues of A are 6i and − 6i. The origin is a stable, but not asymptotically stable, center for the linear part. The theorem does not allow us to draw a definitive conclusion about the almost linear system, which might have a center or spiral point at the origin. Figure 11.26 (a) and (b) shows phase portraits for the almost linear system and its linear part, respectively.

EXAMPLE 11.19

Consider the system




X =

 1

−1 −



 X+

hx
x2 + y2  ky
x2 + y2 

 

y 30 20 10
80


60


40


20

0

20

40

60

80


10
20
30  FIGURE 11.25(a) Phase portrait for

x
= −x + 2y + x sin
y y
= 2x + 3y + 8 sin
x

.

x

438

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y

30 20 10
80


60


40


20

0

20

40

x

60


10
20
30 FIGURE 11.25(b) Phase portrait for the linear part of the system of Figure 11.25(a).

√ 2 in which √ , h and k are constants. The eigenvalues of the matrix of the linear part are  − 1 2 and −  − 1. Consider cases. If 0 <  < 1, then these eigenvalues are pure imaginary. The origin is a center of the linear part but may be a center or spiral point of the almost linear system. If  > 1, then the eigenvalues are real and of opposite sign, so the origin is an unstable saddle point of both the linear part and the original almost linear system. If  = ±1, then A is singular and the system is not almost linear. Figure 11.27 (a) shows a phase portrait for this system with h = 04, k = 07 and  = 13 . Figure 11.27 (b) has  = 2 The next example demonstrates the sensitivity of case (5) of Theorem 11.3. y


30


20


10

0

x

10


10


20


30  FIGURE 11.26(a)

Phase portrait for

x
= 4x + 11y + x sin
y y
= −2x − 4y + sin
y

.

11.5 Almost Linear Systems y 6 4 2


10


5

0

5

x

10


2
4
6 FIGURE 11.26(b)

Phase portrait for the linear part of the system

of Figure 11.26(a). EXAMPLE 11.20

Let  be a real number and consider the system   y + x
x2 + y2   X
= −x + y
x2 + y2  We can write this in the form X
= AX + G as     0 1 x
x2 + y2   X
= X+ y
x2 + y2  −1 0

y 6 4 2
0.8


0.6


0.4


0.2

0

0.2

0.4

0.6

x

0.8


2
4
6
8  FIGURE 11.27(a) Phase portrait for

x
= 13 x − y + 04x
x2 + y2  y
= x − 13 y + 07y
x2 + y2 

.

439

440

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y 6 4 2
0.4


0.2

0

0.2

0.4

x


2
4
6
8 FIGURE 11.27(b)



Phase portrait for




x = 2x − y + 04x
x2 + y2  y
= x − 2y + 07y
x2 + y2 

.

The origin is a critical point of this almost linear system. The eigenvalues of A are i and −i, so the linear part of this system has a center at the origin. This is the case in which Theorem 11.3 does not give a definitive conclusion for the nonlinear system. To analyze the nature of this critical point for the nonlinear system, use polar coordinates r and . Since r 2 = x2 + y 2 then rr
= xx
+ yy




 = x y + x
x2 + y2  + y −x + y
x2 + y2  = 
x2 + y2 
x2 + y2  = r 4  Then dr = r 3  dt This is a separable equation for r, which we solve to get r
t = √

1 k − 2t



in which k is constant determined by initial conditions (a point the trajectory is to pass through). Now consider cases. If  < 0, then r
t = 

1 k + 2  t

→0

as t → . In this case the trajectory approaches the origin in the limit as t → , and
0 0 is asymptotically stable.

11.5 Almost Linear Systems

441

However, watch what happens if  > 0. Say r
0 = , so the trajectory starts at a point at a positive distance  from the origin. Then k = 1/2 and r
t = 

1
1/2  − 2t



In this case, as t increases from 0 and approaches 1/
22 , r
t → . This means that, at finite times, the trajectory is arbitrarily far away from
0 0, hence
0 0 is unstable when  is positive. A phase portrait for  = −02 is given in Figure 11.28 (a), and for  = 02 in Figure 11.28 (b). Figure 11.28 (c) gives a phase portrait for the linear part of this system. y

8 6 4 2
1.0


0.5

0

0.5

x

1.0


2
4  FIGURE 11.28(a) Phase portrait for

with  = −02.

x
= y + x
x2 + y2  y
= −x + y
x2 + y2 

,

y

8 6 4 2
1.0


0.5

0
2
4

FIGURE 11.28(b)

 = 02.

0.5

1.0

x

442

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y 8 6 4 2
8


6


4


2

0
2

2

4

6

8

x


4
6
8

FIGURE 11.28(c)

The linear part
 = 0.

Example 11.20 shows how sensitive an almost linear system can be when the eigenvalues of the linear part are pure imaginary. In this example,  can be chosen arbitrarily small. Still, when  is negative, the origin is asymptotically stable, and when  is positive, regardless of magnitude, the origin becomes unstable. Thus far the discussion has been restricted to nonlinear systems in the special form X
= AX + G, with the origin as the critical point. However, in general a nonlinear system comes in the form X
= F
X, and there may be critical points other than the origin. We will now show how to translate a critical point
x0  y0  to the origin so that X
= F
X translates to a system X
= AX + G. This makes the linear part of the translated system transparent. Further, since Theorem 11.3 is set up to deal with critical points at the origin, we can apply it to X
= AX + G whenever this system is almost linear.   f Thus suppose
x0  y0  is a critical point of X
= F
X, where F = . Assume that f g and g are continuous with continuous first and second partial derivatives at least within some circle about
x0  y0 . By Taylor’s theorem for functions of two variables, we can write, for
x y within some circle about
x0  y0  as f
x y = f
x0  y0  + fx
x0  y0 
x − x0  + fy
x0  y0 
y − y0  + 
x y and g
x y = g
x0  y0  + gx
x0  y0 
x − x0  + gy
x0  y0 
y − y0  + 
x y where lim


xy→
x0 y0 




x y
x − x0

2 +
y − y

0

2

=

lim


xy→
x0 y0 




x y
x − x0 2 +
y − y0 2

= 0

(11.9)

Now
x0  y0  is assumed to be a critical point of X
= F
X, so f
x0  y0  = g
x0  y0  = 0 and these expansions are f
x y = fx
x0  y0 
x − x0  + fy
x0  y0 
y − y0  + 
x y and g
x y = gx
x0  y0 
x − x0  + gy
x0  y0 
y − y0  + 
x y

11.5 Almost Linear Systems Let , X= Then d, X= dt

 

=  =



d
x − x0  dt d
y − y0  dt

f
x y g
x  y



 =

x − x0 y − y0

x
y


443

 

 = X
= F
X



fx
x0  y0 

fy
x0  y0 

gx
x0  y0 

gy
x0  y0 



x − x0 y − y0



 +


x y 
x y



= A
x0 y0  , X + G Because of the condition (11.9), this system is almost linear. Omitting the tilda notation for simplicity, this puts the translated system into the form X
= A
x0 y0  X + G, with the critical point
x0  y0  of X
= F
X translated to the origin as the critical point of the almost linear system X
= A
x0 y0  X + G. Now we can apply the preceding discussion and Theorem 11.3 to the translated system X
= A
x0 y0  X + G at the origin, and hence draw conclusions about behavior of solutions of X
= F
X near
x0  y0 . We use the notation A
x0 y0  for the matrix of the linear part of the translated system for two reasons. First, it reminds us that this is the translated system (since we dropped the , X notation). Second, when we are analyzing several critical points of the same system, this notation reminds us which critical point is under consideration, and clearly distinguishes the linear part associated with one critical point from that associated with another. In carrying out this strategy, it is important to realize that we do not have to explicitly compute 
x y or 
x y, which in some cases would be quite tedious, or not even practical. The point is that we know that the translated system X
= A
x0 y0  X + G is almost linear if F has continuous first and second partial derivatives, a condition that is usually easy to verify.

EXAMPLE 11.21

Consider the system 


X = F
X =

sin
x − x2 + y2



cos

x + y + 1 2 

Here f
x y = sin
x − x2 + y2 and g
x y = cos

x + y + 1/2. This is an almost linear system because f and g are continuous with continuous first and second partial derivatives throughout the plane. For the critical points, solve sin
x − x2 + y2 = 0  cos

x + y + 1  = 0 2 Certainly x = y = n is a solution for every integer n. Every point
n n in the plane is a critical point. There may be other critical points as well, but other solutions of f
x y = g
x y = 0

444

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations

are not obvious. We will need the partial derivatives fx =  cos
x − 2x fy = 2y     sin

x + y + 1  gy = − sin

x + y + 1  2 2 2 2 Now consider a typical critical point
n n. We can translate this point to the origin and write the translated system as X
= A
nn X + G with   fx
n n fy
n n A
nn = gx
n n gy
n n    cos
n − 2n 2n = − 2 sin

2n + 1 2  − 2 sin

2n + 1 2    2n 
−1n − 2n =
−1n+1 2
−1n+1 2 gx = −

and

 G
X =


x y 
x y

 

We need not actually compute 
x y or 
x y. Because the system is almost linear, the qualitative behavior of trajectories of the nonlinear system near
n n is (with exceptions noted in Theorem 11.3) determined by the behavior of trajectories of the linear system X
= A
nn X. We are therefore led to consider the eigenvalues of A
nn , which are 1 1 2 
−1n − n ± 9 − 40n
−1n + 16n2  4 4 We will consider several values for n. For n = 0, the eigenvalues are  and −/2, so the origin is an unstable saddle point of the linear system and also of the nonlinear system. For n = 1, the eigenvalues of A
11 are 1 2 1 9 + 40 + 16 −  −1± 4 4 which are approximately 2 010 1 and −5 580 9. Therefore
1 1 is also an unstable saddle point. For n = 2, the eigenvalues are 1 2 1 9 − 80 + 64  −2± 4 4 which are approximately −1 214 6 + 2 481 2i and −1 214 6 − 2 481 2i. These are complex conjugates with negative real part, so
2 2 is an asymptotically stable spiral point. For n = 3, the eigenvalues are 1 2 1 −  −3± 9 + 120 + 144 4 4 which are approximately −9 959 and 2 388 2. Thus
3 3 is an unstable saddle point. For n = 4, the eigenvalues are 1 1 2  −4± 9 − 160 + 256 4 4 approximately −3 214 6 + 3 140 7i and −3 214 6 − 3 140 7i. We conclude that
4 4 is an asymptotically stable spiral point.

11.5 Almost Linear Systems

445

For n = 5, 1 1 2 9 + 200 + 400 −  −5± 4 4 approximately 2 570 5 and −14 141, so
5 5 is an unstable saddle point. The pattern suggested by the cases n = 2 and n = 4 is broken with n = 6. Now the eigenvalues are 1 1 2 9 − 240 + 576  −6± 4 4 approximately −52146 ± 23606i, so
6 6 is also an unstable spiral point. With n = 7 we get eigenvalues 1 2 1 −  −7± 9 + 280 + 784 4 4 approximately 2 680 2 and −18 251, so
7 7 is an unstable saddle point. The new pattern that seens to be forming is broken with the next case. If n = 8 the eigenvalues are 1 1 2  −8± 9 − 320 + 16
64 4 4 approximately −9 806 9 and −4 622 3, so
8 8 is a stable node. Figures 11.29, 11.30 and 11.31 show phase portraits of this system, focusing on trajectories near selected critical points. The student should experiment with phase portraits near some of the other critical points, for example, those with negative coordinates.

y 3

2

1


3


2


1

0

1

x


1

FIGURE 11.29 Trajectories of the system of Example 11.21 near

the origin.

446

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y

2

1


2

0

2

4

6

8

x


1


2 Trajectories of the system of Example 11.21

FIGURE 11.30

near
1 1.

y

4

3

2

1 1.5

2.0

2.5

3.0

3.5

4.0

4.5

Trajectories of the system of Example 11.21 near

FIGURE 11.31


2 2 and
4 4.

EXAMPLE 11.22 Damped Pendulum

The system for the damped pendulum is x
= y y
= −2 sin
x − !y In matrix form, this system is




X = F
X =

y −2 sin
x − !y

 

x

11.5 Almost Linear Systems

447

Here f
x y = y

and

g
x y = −2 sin
x − !y

The partial derivatives are fx = 0

fy = 1

gx = −2 cos
x

gy = −!

We saw in Example 11.12 that the critical points are
n 0 with n any integer. When n is even, this corresponds to the pendulum bob hanging straight down, and when n is odd, to the bob initially pointing straight up. We will analyze these critical points. Consider first the critical point
0 0. The linear part of the system has matrix   fx
0 0 fy
0 0 A
00 = gx
0 0 gy
0 0   0 1 =  −2 −!   with eigenvalues − 21 ! + 21 ! 2 − 42 and − 21 ! − 21 ! 2 − 42 . Recall that ! = c/mL and 2 = g/L. As we might expect, the relative sizes of the damping force, the mass of the bob and the length of the pendulum will determine the nature of the motion. The following cases occur. (1) If ! 2 − 42 > 0, then the eigenvalues are real, unequal √ and negative, so the origin is an asymptotically stable nodal sink. This happens when c > 2m gL. This gives a measure of how large the damping force must be, compared to the mass of the bob and length of the pendulum, to have trajectories spiralling toward the equilibrium solution
0 0. In this case, after release following a small displacement from the vertical downward position, the bob moves toward this position with decreasing velocity, oscillating back and forth through this position and eventually coming to rest in the limit as t → . Figure 11.32 shows a phase portrait for the pendulum with ! 2 = 08 and  = 044. y

4

2


3


2


1

0

1

2

3


2
4

FIGURE 11.32 Phase portrait for the damped pendulum with ! 2 = 08 and  = 044
! 2 − 42 > 0.

x

448

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y 4

2


3


2


1

0

1

2

x

3


2


4 Damped pendulum with ! 2 = 08 and  = 02
! − 42 = 0. FIGURE 11.33 2

2

(2) If ! 2 − 42 = 0, then the eigenvalues are equal and negative, corresponding to an asymptotically stable proper or improper node of the linear system. This is the case in which Theorem 11.3 does not give a definitive conclusion, and the origin could be an asymptotically stable node or spiral point of the nonlinear pendulum. This case occurs when √ c = 2m gL, a delicate balance between the damping force, mass, and pendulum length. In the case of an asymptotically stable node, the bob, when released, moves with decreasing velocity toward the vertical equilibrium position, approaching it as t → but not oscillating through it. Figure 11.33 gives a phase portrait for this case, in which ! 2 = 08 and 2 = 02. (3) If ! 2 − 42 < 0, then the eigenvalues are complex conjugates with negative real part. Hence the origin is an asymptotically stable spiral√point of both the linear part and the nonlinear pendulum system. This happens when c < 2m gL. Figure 11.34 displays this case, with ! 2 = 06 and 2 = 03. It is routine to check that each critical point
2n 0, in which the first coordinate is an even integer multiple of , has the same characteristics as the origin. Now consider critical points

2n + 1 0, with first coordinate an odd integer multiple of . To be specific, consider
 0. Now the linear part of the system (with
 0 translated to the origin) is  A
0 =

fx
 0 gx
 0

fy
 0 gy
 0



 =

0 1 −2 cos
 −!



 =

0 2

1 −!

 

  The eigenvalues are − 21 ! + 21 ! 2 + 42 and − 21 ! − 21 ! 2 + 42 . These are real and of opposite sign, so
 0 is an unstable saddle point. The other critical points

2n + 1 0 exhibit the same behavior. This is what we would expect of a pendulum in which the bob is initially in the vertical upward position, since arbitrarily small displacements will result in the bob moving away from this position, and it will never return to it. The analysis is the same for each critical point

2n + 1 0.

11.5 Almost Linear Systems

449

y

2 1


1.5


1.0


0.5

0

0.5

1.0

1.5

x


1
2

FIGURE 11.34 Damped pendulum with ! 2 = 06 and

2 = 03
! 2 − 42 < 0.

EXAMPLE 11.23 Nonlinear Damped Spring

The nonlinear damped spring equation is x
= y y
= −

 k c x + x3 − y m m m

This is 


X = F
X =

y −
k/mx +
/mx3 −
c/my

 

Here

f
x y = y and g
x y = −

k  c x + x3 − y m m m

√ √ There are three critical points,
0 0,
k/ 0 and
− k/ 0. The partial derivatives are

fx = 0 fy = 1 gx = −

k c  + 3 x 2  gy = −  m m m

450

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations

First consider the behavior of trajectories near the origin. The linear part of the system has matrix   0 1 A
00 =  −k/m −c/m   √ √ with eigenvalues
1/2m −c + c2 − 4mk and
1/2m −c − c2 − 4mk . This yields three cases, depending, as we might expect, on the relative magnitudes of the mass, damping constant and spring constant. 1. If c2 − 4mk > 0, then A
00 has real, distinct, negative eigenvalues, so the origin is an asymptotically stable nodal sink. Small disturbances from the equilibrium position result in a motion that dies out with time, with the mass approaching the equilibrium position. 2. If c2 − 4mk = 0, then A
00 has equal real, negative eigenvalues, so the origin is an asymptotically stable proper or improper node of the linear system. Hence the origin is an asymptotically stable node or spiral point of the nonlinear system. 3. If c2 − 4mk < 0, then A
00 has complex conjugate eigenvalues with negative real part, and the origin is an asymptotically stable spiral point. Figure 11.35 shows a phase portrait √ for case (3), with c = 2, k = 5,  = 1 and m = 3. Next, consider the critical point
k/ 0. Now the linear part of the system obtained by translating this point to the origin is   0 1 √  A
k/0 = k/m −c/m

y

2

1


2


1

0

1

x


1

Nonlinear spring system with c = 2, k = 5,  = 1, and m = 3
c2 − 4mk < 0.

FIGURE 11.35

  √ √ with eigenvalues
1/2m −c + c2 + 4mk and
1/2m −c − c2 + 4mk . The first eigen√ value is positive and the second negative, √ so
k/ 0 is an unstable saddle point. A similar analysis holds for the critical point
− k/ 0

11.6 Lyapunov’s Stability Criteria

SECTION 11.5

PROBLEMS

In each of Problems 1 through 10, (a) show that the system is almost linear, (b) determine the critical points, (c) use Theorem 11.3 to analyze the nature of the critical point, or state why no conclusion can be drawn, and (d) generate a phase portrait for the system.


is possible in this case in general by considering the following two systems:   x
= y − x x2 + y2  y
= −x − y x2 + y2 and




1. x = x − y + x  y = x + 2y 2

451

2. x
= x + 3y − x2 sin
y y
= 2x + y − xy2

  x
= y + x x2 + y2  y
= −x + y x2 + y2 

(a) Show that the origin is a center for the associated linear system of both systems. (b) Show that each system is almost linear.

3. x
= −2x + 2y y
= x + 4y + y2

(c) Introduce polar coordinates, with x = r cos
 and y = r sin
 and use the chain rule to obtain

4. x
= −2x − 3y − y2  y
= x + 4y

x
=

5. x
= 3x + 12y y
= −x − 3y + x3





6. x = 2x − 4y + 3xy y = x + y + x

and 2

7. x
= −3x − 4y + x2 − y2  y
= x + y 8. x
= −3x − 4y y
= −x + y − x2 y 9. x
= −2x − y + y2  y
= −4x + y 10. x
= 2x − y − x3 sin
x y
= −2x + y + xy2 11. Theorem 11.3 is inconclusive in the case that the critical point of an almost linear system is a center of the associated linear system. Verify that no conclusion

11.6

dx dr = cos
r

t dr dt

dy dr = sin
r

t dr dt Use these to evaluate xx
+ yy
in terms of r and r
, where r
= dr/dt. Thus convert each system to a system in terms of r
t and 
t y
=

(d) Use the polar coordinate version of the first system to obtain a separable differential equation for r
t. Conclude from this that r

t < 0 for all t. Solve for r
t and show that r
t → 0 as t → . Thus conclude that for the first system the origin is asymptotically stable. (e) Follow the procedure of (d), using the second system. However, now find that r

t > 0 for all t. Solve for r
t with the initial condition r
t0  = r0 . Show that r
t → as t → t0 + 1/r from the left. Conclude that the origin is unstable for the second system.

Lyapunov’s Stability Criteria There is a subtle criterion for stability due to the Russian engineer and mathematician Alexander M. Lyapunov (1857-1918). Suppose X
= F
X is a 2 × 2 autonomous system of first order differential equations (not necessarily almost linear), and that
0 0 is an isolated critical point. Lyapunov’s insight was this. Suppose there is a function, commonly denoted V , such that closed curves V
x y = c enclose the origin. Further, if the constants are chosen smaller, say 0 < k < c, then the curve V
x y = k lies within the region enclosed by the curve V
x y = c (Figure 11.36 (a)). So far this has nothing to do with the system of differential equations. However, suppose it also happens that, if a trajectory intersects the curve V
x y = c at some time, which we can take to be time zero, then it cannot escape from the region bounded by this curve, but must for all later times remain within this region (Figure 11.36 (b)). This would force trajectories starting out near the origin (meaning within V
x y = c) to forever lie at least this

452

CHAPTER 11

Qualitative Methods and Systems of Nonlinear Differential Equations y V(x, y) = 4 V(x, y) = 2 x V(x, y) = 1 FIGURE 11.36(a)

Closed curves contracting

about the origin.

y

V(x, y) = c

x

(x(t), y(t))

(x(0), y(0))

FIGURE 11.36(b) Trajectories entering shrinking regions about the origin.

close to the origin. But this would imply, by choosing c successively smaller, that the origin is a stable critical point! If in addition trajectories starting at a point on V
x y = c point into the region bounded by this curve, then we can further conclude that the trajectories are approaching the origin, hence that the origin is asymptotically stable. This is the intuition behind an approach to determining whether a critical point is stable or asymptotically stable. We will now develop the vocabulary which will allow us to give substance to this approach. First, we will distinguish certain functions that have been found to serve the role of V in this discussion. If r > 0, let Nr consist of all
x y within distance r from the origin. Thus,
x y is in Nr exactly when x2 + y 2 < r 2  This set is called the r-neighborhood of the origin, or, if we need no explicit reference to r, just a neighborhood of the origin.

11.6 Lyapunov’s Stability Criteria

DEFINITION 11.5

453

Positive Definite, Semidefinite

Let V
x y be defined for all
x y in some neighborhood Nr of the origin. Suppose V is continuous with continuous first partial derivatives. Then 1. V is positive definite on Nr if V
0 0 = 0 and V
x y > 0 for all other points of Nr . 2. V is positive semidefinite on Nr if V
0 0 = 0 and V
x y ≥ 0 for all points of Nr . 3. V is negative definite on Nr if V
0 0 = 0 and V
x y < 0 for all other points of Nr . 4. V is negative semidefinite on Nr if V
0 0 = 0 and V
x y ≤ 0 for all points of Nr .

For example, V
x y = x2 + 3xy + 9y2 is positive definite on Nr for any positive r, and −3x2 + 4xy − 5y2 is negative definite on any Nr . The function
x − y4 is positive semidefinite, being nonnegative but vanishing on the line y = x. The following lemma is useful in producing examples of positive definite and negative definite functions.

LEMMA 11.2

Let V
x y = ax2 + bxy + cy2 . Then V is positive definite (on any Nr ) if and only if a > 0 and 4ac − b2 > 0 V is negative definite (on any Nr ) if and only if a < 0 and 4ac − b2 > 0 Certainly V is continuous with continuous partial derivatives in the entire x y plane. Further, V
0 0 = 0, and this is the only point at which V
x y vanishes. Now recall the second derivative test for extrema of a function of two variables. First,

Proof

Vx
0 0 = Vy
0 0 = 0 so the origin is a candidate for a maximum or minimum of V . For a maximum or minimum, we need Vxx
0 0Vyy
0 0 − Vxy
0 02 > 0 But this condition is the same as
2a
2c − b2 > 0 or 4ac − b2 > 0 Satisfaction of this inequality requires that a and c have the same sign.

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When a > 0, then Vxx
0 0 > 0 and the origin is a point where V has a minimum. In this event V
x y > V
0 0 = 0 for all
x y other than
0 0. Now V is positive definite. When a < 0 then Vxx
0 0 < 0 and the origin is a point where V has a maximum. Now V
x y < 0 for all
x y other than
0 0, and V is negative definite. If x = x
t y = y
t defines a trajectory of X
= F
X, and V is a differentiable function of two variables, then V
x
t y
t is a differentiable function of t along this trajectory. We will denote the derivative of V
x
t y
t with respect to t as V˙
x y, or just V˙ . By the chain rule, V˙
x y = Vx
x
t y
tx

t + Vy
x
t y
ty

t or, more succinctly, V˙ = Vx x
+ Vy y
 This is called the derivative of V along the trajectory, or the orbital derivative of V . The following two theorems show how these ideas about positive and negative definite functions relate to Lyapunov’s approach to stable and asymptotically stable critical points. The criteria given in the first theorem constitute Lyapunov’s direct method for determining the stability or asymptotic stability of a critical point. THEOREM 11.4

Lyapunov’s Direct Method for Stability

Let
0 0 be an isolated critical point of the autonomous 2 × 2 system X
= F
X. 1. If a positive definite function V can be found for some neighborhood Nr of the origin, such that V˙ is negative semidefinite on Nr , then the origin is stable. 2. If a positive definite function V can be found for some neighborhood Nr of the origin, such that V˙ is negative definite on Nr , then the origin is asymptotically stable. On the other side of the issue, Lyapunov’s second theorem gives a test to determine that a critical point is unstable. THEOREM 11.5

Lyapunov’s Direct Method for Instability

Let
0 0 be an isolated critical point of the autonomous 2 × 2 system X
= F
X. Let V be continuous with continuous first partial derivatives in some neighborhood of the origin, and let V
0 0 = 0. 1. Suppose R > 0, and that in every neighborhood Nr , with 0 < r ≤ R, there is a point at which V
x y is positive. Suppose V˙ is positive definite in NR . Then
0 0 is unstable. 2. Suppose R > 0, and that in every neighborhood Nr , with 0 < r ≤ R, there is a point at which V
x y is negative. Suppose V˙ is negative definite in NR . Then
0 0 is unstable. Any function V playing the role cited in these theorems is called a Lyapunov function. Theorems 11.4 and 11.5 give no suggestion at all as to how a Lyapunov function might be produced, and in attempting to apply them this is the difficult part. Lemma 11.2 is sometimes useful in providing candidates, but, as might be expected, if the differential equation is complicated the task of finding a Lyapunov function might be insurmountable. In spite of this potential difficulty, Lyapunov’s theorems are useful because they do not require solving the system, nor do they require that the system be almost linear.

11.6 Lyapunov’s Stability Criteria

455

Adding to the mystique of the theorem is the nonobvious connection between V , V˙ , and stability characteristics of the critical point. We will give a plausibility argument intended to clarify this connection. Consider Figure 11.37, which shows a typical curve V
x y = c about the origin. Call this curve $. Here is how V˙ enters the picture. At any point P 
a b on this curve, the vector N = Vx i + Vy j is normal (perpendicular) to $, by which we mean that it is normal to the tangent to $ at this point. In addition, consider a trajectory x = 
t, y = 
t passing through P at time t = 0, also shown in Figure 11.37. Thus, 
0 = a and 
0 = b. The vector T = 

0i + 

0j is tangent to this trajectory (not to $) at
a b. Now, V˙
a b = Vx
a b

0 + Vy
a b

0 = N · T the dot product of the normal to $ and the tangent to the trajectory at
a b. Since the dot product of two vectors is equal to the product of their lengths and the cosine of the angle between them, we obtain V˙
a b = N T cos
 with  the angle between T and N. Now look at conclusions (1) and (2) of the first Lyapunov theorem. If V˙ is negative semidefinite, then V˙
a b ≤ 0, so cos
 ≤ 0. Then /2 ≤  ≤ 3/2. This means that the trajectory at this point is moving at this point either into the region enclosed by $, or perhaps in the same direction as the tangent to $. The effect of this is that the trajectory cannot move away from the region enclosed by $. The trajectory cannot escape from this region, and so the origin is stable. If V˙ is negative definite, then cos
 < 0, so /2 <  < 3/2 and now the trajectory actually moves into the region enclosed by $, and cannot simply trace out a path around the origin. In this case the origin is asymptotically stable. We leave it for the student to make a similar geometric argument in support of the second Lyapunov theorem.

y

Γ: V(x, y) = c

x P N

θ

T Tangent to Γ at P

FIGURE 11.37 Rationale for Lyapunov’s direct method.

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EXAMPLE 11.24

Consider the nonlinear system




X =

−x3 −4x2 y

 

The origin is an isolated critical point. We will attempt to construct a Lyapunov function of the form V
x y = ax2 + bxy + cy2 that will tell us whether the origin is stable or unstable. We may not succeed in this, but it is a good first attempt because at least we know conditions on the coefficients of V to make this function positive definite or negative definite. The key lies in V˙ , so compute V˙ = Vx x
+ Vy y
=
2ax + by
−x3  +
bx + 2cy
−4x2 y = −2ax4 − bx3 y − 4bx3 y − 8cx2 y2  Now observe that the −2ax4 and −8cx2 y2 terms will be nonpositive if a and c are positive. The x3 y term will vary in sign, but we can make bx3 y vanish by choosing b = 0. We can choose a and c as any positive numbers, say a = c = 1. Then V
x y = x2 + y2 is positive definite in any neighborhood of the origin , and V˙ = −2x4 − 8x2 y2 is negative semidefinite in any neighborhood of the origin. By Lyapunov’s direct method (Theorem 11.4), the origin is stable. A phase portrait for this system is shown in Figure 11.38. We can draw no conclusion about asymptotic stability of the origin in the last example. If we had been able to find a Lyapunov function V so that V˙ was negative definite (instead of

y 2

1


1.0


0.5

0

0.5

1.0


1


2  FIGURE 11.38

Phase portrait for

x
= −x3 y
= −4x2 y

.

x

11.6 Lyapunov’s Stability Criteria

457

negative semidefinite), then we could have concluded that the origin was asymptotically stable. But we cannot be sure, from the work done, whether there is no such function, or whether we simply did not find one.

EXAMPLE 11.25

It is instructive to consider Lyapunov’s theorems as they relate to a simple physical problem, the undamped pendulum (put c = 0). The system is x
= y g sin
x L Recall that x = , the displacement angle of the bob from the downward vertical rest position. We have already characterized the critical points of this problem (with damping). However, this example makes an important point. In problems that are drawn from a physical setting, the total energy of the system can often serve as a Lyapunov function. This is a useful observation, since the search for a Lyapunov function constitutes the primary issue in attempting to apply Lyapunov’s theorems. Thus compute the total energy V of the pendulum. The kinetic energy is y=−

1 2 mL2


t  2 which in the variables of the system is 1 mL2 y2  2 The potential energy is the work done in lifting the bob above the lowest position. From Figure 11.1, this is mgL
1 − cos
 or mgL
1 − cos
x The total energy is therefore given by 1 V
x y = mgL
1 − cos
x + mL2 y2  2 Clearly V
0 0 = 0. The rest position of the pendulum (pendulum arm vertical with bob at the low point) has zero energy. Next, compute V˙
x
t y
t = Vx x

t + Vy y

t = mgL sin
xx

t + mL2 yy

t Along any trajectory of the system, x
= y and y
= −
g/L sin
x, so the orbital derivative is  g V˙
x
t y
t = mgL sin
xy + mL2 y − sin
x = 0 L This corresponds to the fact that, in a conservative physical setting, the total energy is a constant of the motion. Now V is positive definite. This is expected because the energy should be a minimum in the rest position, where V
0 0 = 0. Further, V˙ is negative semidefinite. Therefore the origin is stable.

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As expected, Lyapunov’s theorem did not tell us anything new about the pendulum, which we had already analyzed by other means. However, this example does provide some insight into a line of thought that could have motivated Lyapunov. In any conservative physical system, the total energy must be a constant of the motion, and we expect that any position with the system at rest should be stable if the potential energy is a minimum, and unstable if it is not. This suggests looking at the total energy as a candidate for a Lyapunov function V . In particular, for many mechanical systems the kinetic energy is a quadratic form. One then checks the orbital derivative V˙ to see if it is negative definite or semidefinite in a neighborhood of the point of interest.

EXAMPLE 11.26

Consider the system

 X
=

x3 − xy2 y3 + 6x2 y

 

The origin is an isolated critical point. We do not know whether it is stable, asymptotically stable, or unstable, so we will begin by trying to construct a Lyapunov theorem that fits either of the Lyapunov theorems. Attempt a Lyapunov function V
x y = ax2 + bxy + cx2  We know how to choose the coefficients to make this positive definite. The question is what happens with the orbital derivative. Compute V˙ = Vx x
+ Vy y
=
2ax + by
x3 − xy2  +
bx + 2cy
y3 + 6x2 y = 2ax4 + 2cy4 +
12c − 2ax2 y2 + 7bx3 y This looks promising because the first three terms can be made strictly positive for
x y =
0 0. Thus choose b = 0 and a = c = 1 to get V˙
x y = 2x4 + 2y4 + 10x2 y2 > 0 for
x y =
0 0. With this choice of the coefficients, V
x y = x2 + y2 is positive definite on any neighborhood of the origin, and V˙ is also positive definite. By Lyapunov’s second theorem (Theorem 11.5), the origin is unstable. Figure 11.39 shows a phase portrait for this system.

EXAMPLE 11.27

A nonlinear oscillator with linear damping can be modeled by the differential equation z

+ z
+ z + z2 + !z3 = 0 in which  and ! are positive and  = 0. To convert this to a system, let z = x and z
= y to obtain x
= y y
= −y − x − x2 − !x3 

11.6 Lyapunov’s Stability Criteria

459

y 2 1.5 1 0.5
1.0


0.5

0

0.5

1.0

x


0.5
1  FIGURE 11.39 Phase portrait for

x
= x3 − xy2 y
= y3 + 6x2 y

.

This is the system  X
=

y −y − x − x2 − !x3

 

We can construct a Lyapunov function by a clever observation. Let 1 1 1 1 V
x y = y2 + x2 + x3 + !x4  2 2 3 4 Then V˙ =
x + x2 + !x3 y +
y
−y − x − x2 − !x3  = −y2 . Since  > 0, V˙ is certainly negative semidefinite in any neighborhood of the origin. It may not be obvious whether V is positive definite in any neighborhood of the origin. Certainly the term y2 /2 in V is nonnegative. The other terms are 1 2 1 3 1 4 x + x + !x  2 3 4 which we can write as

 x

2

 1 1 2 1 + !x + x  2 4 3

Since g
x = 21 + 41 !x2 + 13 x is continuous for all x, and g
0 = 1/2, there is an interval
−h h about the origin such that g
x > 0 for − h < x < h

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Then, in Nh , V
x y ≥ 0, and V
x y > 0 if
x y =
0 0. Therefore V is positive definite in this neighborhood. We now have V positive definite and V˙ negative semidefinite in Nh , hence the origin is stable. Figure 11.40 shows a phase portrait for the case  = 1,  = 41 and ! = 16 . It is instructive to try different values of  and ! to get some idea of effect of the nonlinear terms on the trajectories.

y 40 20


15


10


5

0

5

10

x


20
40
60

FIGURE 11.40

Nonlinear oscillator with  = 1,  = 41 , ! = 16 .

EXAMPLE 11.28

Consider the system x
= f
ty + g
tx
x2 + y2  y
= −f
tx + g
ty
x2 + y2  Assume that the origin is an isolated critical point. If we attempt a Lyapunov function V
x y = ax2 + bxy + cy2 , then V˙ =
2ax + by f
ty + g
tx
x2 + y2  +
bx + 2cy −f
tx + g
ty
x2 + y2  =
2a − 2cf
txy + 2ax2 g
t
x2 + y2  + 2cy2 g
t
x2 + y2  + bf
ty2 + 2bg
txy
x2 + y2  − bf
tx2  We can eliminate three terms in the orbital derivative by trying b = 0. The term
2a − 2c f
txy vanishes if we choose a = c. To have V positive definite we need a and c positive, so let a = c = 1. Then V
x y = x2 + y2  which is positive definite in any neighborhood of the origin, and V˙ = 2
x2 + y2 2 g
t

11.7 Limit Cycles and Periodic Solutions

461

This is negative definite if g
t < 0 for all t ≥ 0. In this case the origin is asymptotically stable. If g
t ≤ 0 for t ≥ 0, then V is negative semidefinite and then the origin is stable. If g
t > 0 for t ≥ 0, then the origin is unstable.

SECTION 11.6

PROBLEMS

In each of Problems 1 through 8, use Lyapunov’s theorem to determine whether the origin is stable asymptotically stable, or unstable.

4. x
= −x2 y2  y
= x2

1. x
= −2xy2  y
= −x2 y

6. x
= x5
1 + y2  y
= x2 y + y3

2. x
= −x cos2
y y
=
6 − xy2

7. x
= x3
1 + y y
= y3
4 + x2 

3. x
= −2x y
= −3y3

8. x
= x3 cot 2
y y
= y3
2 + x4 

11.7

5. x
= xy2  y
= y3

Limit Cycles and Periodic Solutions Nonlinear systems of differential equations can give rise to curves called limit cycles which have particularly interesting properties. To see how a limit cycle occurs naturally in a physical setting, draw a circle C of radius R on pavement, with R exceeding the length L between the points where the front and rear wheel of a bicycle touch the ground. Now grab the handlebars and push the bicycle so that its front wheel moves around C. What path does the rear wheel follow? If you tie a marker to the rear wheel so that it traces out the rear wheel’s path as the front wheel moves along C, you find that this path does not approach a particular point. Instead, as the front wheel continues its path around√C, the rear wheel asymptotically approaches a circle K concentric with C and having radius R2 − L2 . If the rear wheel begins outside C, it will spiral inward toward K, while if it begins inside C , it will work its way outward toward K. If the rear wheel begins on K, it will remain on K. This inner circle K has two properties in common with a stable critical point. Trajectories beginning near K move toward it, and if a trajectory begins on K, it remains there. However, K is not a point, but is instead a closed curve. K is a limit cycle of this motion.

DEFINITION 11.6

Limit Cycle

A limit cycle of a 2 × 2 system X
= F
X is a closed trajectory K having the property that there are trajectories x = 
t, y = 
t of the system such that

t 
t spirals toward K in the limit as t → .

We have already pointed out the analogy between a limit cycle and a critical point. This analogy can be pushed further by defining a concept of stability and asymptotic stability for limit cycles that is modeled after stability and asymptotic stability for critical points.

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DEFINITION 11.7

Let K be a limit cycle of X
= F
X. Then, 1. K is stable if trajectories starting within a certain distance of K must remain within a fixed distance of K. 2. K is asymptotically stable if every trajectory that starts sufficiently close to K spirals toward K as t → . 3. K is semistable if every trajectory starting on one side of K spirals toward K as t → , while there are trajectories starting on the other side of K that spiral away from K as t → . 4. K is unstable if there are trajectories starting on both sides of K that spiral away from K as t → .

Keep in mind that a closed trajectory of X
= F
X represents a periodic solution. We have seen periodic solutions previously with centers, which are critical points about which trajectories form closed curves. A limit cycle is therefore a periodic solution, toward which other solutions approach spirally. Often we are interested in whether a system X
= F
X has a periodic solution, and we will shortly develop some tests to tell whether a system has such a solution, or sometimes to tell that it does not.

EXAMPLE 11.29 Limit Cycle

Consider the almost linear system 
X =

1 −1

1 1



 X+

  −xx2 + y2  −y x2 + y2

(11.10)


0 0 is the only critical point of this system. The eigenvalues of   1 1 −1 1 are 1 ± i, so the origin is an unstable spiral point of the system X
= AX + G and also of the linear system X
= AX. Figure 11.41 shows a phase portrait of the linear system. Up to this point, whenever we have seen trajectories spiraling outward, they have grown without bound. We will now see that this does not happen with this current system. This will be transparent if we convert the system to polar coordinates. Since r 2 = x2 + y 2  then rr
= xx
+ yy
      = x x + y − x x2 + y2 + y −x + y − y x2 + y2  = x2 + y2 −
x2 + y2  x2 + y2 = r 2 − r 3 

11.7 Limit Cycles and Periodic Solutions

463

y 10

5


10


5

0

5

10

x


5


10 FIGURE 11.41 Unstable spiral point for x
= x + y, y
= −x + y.

Then r
= r − r 2 = r
1 − r This tells us that, if 0 < r < 1, then r
> 0 so the distance between a trajectory and the origin is increasing. Trajectories inside the unit circle are moving outward. But if r > 1, then r
< 0, so the distance between a trajectory and the origin is decreasing. Trajectories outside the unit circle are moving inward. This does not yet tell us in detail how the trajectories are moving outward or inward. To determine this, we need to bring the polar angle  into consideration. Differentiate x = r cos
 and y = r sin
 with respect to t: x
= r
cos
 − r sin

 y
= r
sin
 + r cos

 Now observe that x
y − xy
= r sin
 r
cos
 − r sin

 − r cos
 r
sin
 + r cos



 = −r 2 cos2
 + sin2
 
= −r 2 
 But from the system X
= AX + G we have       x
y − xy
= y x + y − x x2 + y2 − x −x + y − y x2 + y2 = x2 + y 2 = r 2  Therefore, r 2 = −r 2 


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from which we conclude that 
= −1 We now have an uncoupled system of differential equations for r and : r
= r
1 − r


= −1

This equation for r is separable. For the trajectory through
r0  0 , solve these equations subject to the initial conditions r
0 = r0 


0 = 0 

We get r=

1   = 0 − t 1 −

r0 − 1/r0 e−t

These explicit solutions enable us to conclude the following. If r0 = 1, then
r0  0  is on the unit circle. But then r
t = 1 for all t, hence a trajectory that starts on the unit circle remains there for all times. Further, 

t = −1, so the point
r
t 
t moves clockwise around this circle as t increases. If 0 < r0 < 1, then
r0  0  is within the disk bounded by the unit circle. Now r
t =

1 0. Therefore a trajectory starting at a point inside the unit disk remains there forever. Further, r
t → 1 as t → , so this trajectory approaches the unit circle from within. Finally, if r0 > 1, then
r0  0  is outside the unit circle. But now r
t > 1 for all t, so a trajectory starting outside the unit circle remains outside for all time. However, it is still true that r
t → 1 as t → , so this trajectory approaches the unit circle from without. In sum, trajectories tend in the limit to wrap around the unit circle, either from within or from without, depending on where they start.The unit circle is a an asymptotically stable limit cycle of X
= F
X. A phase portrait is shown in Figure 11.42. Here is an example of a system having infinitely many asymptotically stable limit cycles. EXAMPLE 11.30

The system  X
=

0 −1

1 0



⎛ X+⎝

x sin y sin



x2 + y 2

⎞

⎠  x2 + y 2

has particularly interesting limit cycles, as can be seen in Figure 11.43. These occur as concentric circles about the origin. Trajectories originating within the innermost circle spiral toward this circle, as do trajectories beginning between this first circle and the second one. Trajectories beginning between the second and third circles spiral toward the third circle, as do trajectories originating between the third and fourth circles. This pattern continues throughout the plane. We will now develop some facts about closed trajectories (periodic solutions) and limit cycles. For the remainder of this section,   f
x y
X = F
X = g
x y is a 2 × 2 autonomous system.

11.7 Limit Cycles and Periodic Solutions y 4

2


3


2

0

2

3

x


2


4 FIGURE 11.42 Limit cycle of   x
= x + y − x x2 + y 2 .  y
= −x + y − y x2 + y2

y

10 5


10


5

0

5


5
10

FIGURE 11.43 Asymptotically stable limit cycles of   x
= y + x sin
x2 + y2  .  y
= −x + y sin
x2 + y2 

10

x

465

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The first result states that, under commonly encountered conditions, a closed trajectory of X
= F
X must always enclose a critical point. THEOREM 11.6

Enclosure of Critical Points

Let f and g be continuous with continuous first partial derivatives in a region of the plane containing a closed trajectory K of X
= F
X. Then K must enclose at least one critical point of X
= F
X. This kind of result can sometimes be used to tell that certain regions of the plane cannot contain closed trajectories of a system of differential equations. For example, suppose the origin is the only critical point of X
= F
X. Then we automatically know that there can be no closed trajectory in, for example, one of the quadrants, because such a closed trajectory could not enclose the origin, contradicting the fact that it must enclose a critical point. Bendixson’s theorem, which follows, gives conditions under which X
= F
X has no closed trajectory in a part of the plane. A region of the plane is called simply connected if it contains all the points enclosed by any closed curve in the region. For example, the region bounded by the unit circle is simply connected. But the shaded region shown in Figure 11.44 between the curves C and K is not simply connected, because C encloses points not in the region. A simply connected region can have no “holes” in it, because then a closed curve wrapping around a hole encloses points not in the region. y K

C x

FIGURE 11.44 Non-simply connected region.

THEOREM 11.7

Bendixson

Let f and g be continuous with continuous first partial derivatives in a simply connected region R of the plane. Suppose fx + gy has the same sign throughout points of R, either positive or negative. Then X
= F
X has no closed trajectory in R. Proof Suppose R contains a closed trajectory C representing the periodic solution x = 
t, y = 
t. Suppose

t 
t traverses this curve exactly once as t varies from a to b, and let D be the region enclosed by C. Evaluate the line integral:
b −g
x ydx + f
x ydy =

−g

t 
t

t + f

t 
t

t dt C

a

So far we have not used the fact that x = 
t, y = 
t is a solution of X
= F
X. Using this, we have 

t = x
= f
x y = f

t 
t

11.7 Limit Cycles and Periodic Solutions

467

and 

t = y
= g
x y = g

t 
t Then −g

t 
t

t + f

t 
t

t = −g

t 
tf

t 
t + f

t 
tg

t 
t = 0 Therefore

C

But by Green’s theorem, C

−g
x ydx + f
x ydy = 0

−g
x ydx + f
x ydy =



 D

 fx + gy dxdy

and this integral cannot be zero because the integrand is continuous and of the same sign throughout D. This contradiction implies that no such closed trajectory C can exist within the region R.

EXAMPLE 11.31

Consider the system




X =

3x + 4y + x3 5x − 2y + y3

 

Here f and g are continuous, with continuous first partial derivatives, throughout the plane. Further, fx + gy = 3 + 3x2 − 2 + 3y2 > 0 for all
x y. This system has no closed trajectory, hence no periodic solution. The last two theorems have been negative, in the sense of providing criteria for X
= F
X to have no periodic solution in some part of the plane. The next theorem, a major result credited dually to Henri Poincaré and Ivar Bendixson, gives a condition under which X
= F
X has a periodic solution. THEOREM 11.8

Poincaré-Bendixson

Let f and g be continuous with continuous first partial derivatives in a closed, bounded region R of the plane that contains no critical point of X
= F
X. Let C be a trajectory of X
= F
X that is in R for t ≥ t0 . Then, C must be a periodic solution (closed trajectory), or else C spirals toward a closed trajectory as t → . In either case, as long as a trajectory enters R at some time, and R contains no critical point, then X
= F
X has a periodic solution, namely this trajectory itself, or, if not, a closed trajectory approached spirally by this trajectory. On the face of it, this result may appear to contradict the conclusion of Theorem 11.6, since any periodic trajectory should enclose a critical point. However, this critical point need not be in the region R of the theorem. To illustrate, consider again the system (11.10). Let R be the region between the concentric circles r = 1 and r = 3 shown in Figure 11.45. The only critical point of X
= F
X is the origin, which is not in R. If we choose any trajectory

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r3 x

FIGURE 11.45 Region between the circles r = 21 and r = 3, with trajectories approaching the limit cycle r = 1.

beginning at a point inside the unit circle, then, as we have seen, this trajectory approaches the unit circle, hence eventually enters R. The Poincaré -Bendixson theorem would allow us to assert just from this that R contains a periodic solution of the system. We will conclude this section with Lienard’s theorem, which gives conditions sufficient for a system to have a limit cycle. THEOREM 11.9

Lienard

Let p and q be continuous and have continuous derivatives on the entire real line. Suppose: 1. q
x = −q
−x for all x. 2. q
x > 0 for all x > 0. 3. p
x = p
−x for all x. Suppose also the equation F
x = 0 has exactly one positive root, where
x F
x = p
d 0

If this root is denoted !, suppose F
x < 0 for 0 < x < !, and F
x is positive and nondecreasing for x > !. Then, the system   y X
= −p
xy − q
x has a unique limit cycle enclosing the origin. Further, this limit cycle is asymptotically stable. Under the conditions of the theorem, the system has exactly one periodic solution, and every other trajectory spirals toward this closed curve as t → . As an illustration, we will use Lienard’s theorem to analyze the van der Pol equation. EXAMPLE 11.32 van der Pol Equation

The second-order differential equation z

+ 
z2 − 1z
+ z = 0

11.7 Limit Cycles and Periodic Solutions

469

in which  is a positive constant, is called van der Pol’s equation. It was derived by the Dutch engineer Balthazar van der Pol in the 1920’s in his studies of vacuum tubes. It was of great interest to know whether this equation has periodic solutions, a question to which the answer is not obvious. First write van der Pol’s equation as a system. Let x = z and y = z
to get x
= y y
= −
x2 − 1 − x This system has exactly one critical point, the origin. Further, this system matches the one in Lienard’s theorem if we let p
x = 
x2 − 1 and q
x = x. Now q
−x = −q
x, q
x > 0 for x > 0, and p
−x = p
x, as required. Next, let  
x 1 2 F
x = x −1  p
d = x 3 0 √ √ √ F has exactly one positive zero, ! = 3. For 0 < x < 3, F
x < 0. Further, for x > 3, F
x is positive and increasing (hence nondecreasing). By Lienard’s theorem, the van der Pol equation has a unique limit cycle (hence periodic solution) enclosing the origin. This limit cycle is asymptotically stable, so trajectories beginning at points not on this closed trajectory spiral toward this limit cycle. Figures 11.46 through 11.49 show phase portraits for van der Pol’s equation, for various choices of . y

y

6

6

4

4

2


6


4


2

2

2

4

6

x


6


4


2

2 4


2


4


2


4


6


6

FIGURE 11.46 Phase portrait for van der Pol’s equation for  = 02.

FIGURE 11.47 Phase portrait for van der Pol’s equation for  = 05.

6

x

CHAPTER 11

470

Qualitative Methods and Systems of Nonlinear Differential Equations x2 y 3 4

2 2 1


4


2

0

2

4

x


2


1

0

1


1
2


2


3


4 Phase portrait for van der Pol’s equation

FIGURE 11.48

for  = 1.

SECTION 11.7

PROBLEMS

In each of Problems 1 through 4, use Bendixson’s theorem to show that the system has no closed trajectory. Generate a phase portrait for the system.





1. x = −2x − y + x  y = 10x + 5y + x y − 2y sin
x 3

2

2. x
= −x − 3y + e2x  y
= x + 2y + cos
y 3. x
= 3x − 7y + sinh
x y
= −4y + 5e3y 4. x
= y y
= −x + y
9 − x2 − y2 , for
x y in the elliptical region bounded by the graph of x2 + 9y2 = 9 5. Recall that with the transformation from rectangular to polar coordinates, we obtain x

dx dy dr +y =r dt dt dt

and y

FIGURE 11.49 Phase portrait for van der Pol’s equation for  = 3.

dy d dx −x = −r 2  dt dt dt

Use these equations to show that the system x
= y + 

x x2 + y 2

 f
x2 + y2 

 y f
x2 + y2  y
= −x +  2 2 x +y has closed trajectories associated with zeros of the function f (a given continuous function of one variable). If the closed trajectory is a limit cycle, what is its direction of orientation? In each of Problems 6 through 9, use a conversion to polar coordinates (see Problem 5) to find all of the closed trajectories of the system, and determine which of these are limit cycles. Also classify the stability of each limit cycle. Generate a phase portrait for each system and attempt to identify the limit cycles.

2

x1

11.7 Limit Cycles and Periodic Solutions   6. x
= 4y + x sin
x2 + y2  y
= −x + y sin
x2 + y2 


18. x
= y y
= −x + y − x2 y

8. x
= x
1 − x2 − y2  y
= y
1 − x2 − y2 

19. x
= x − 5y + y3  y
= x − y + y3 + 7y5

9. x
= y + x
1 − x2 − y2 
4 − x2 − y2 
9 − x2 − y2  y
= −x + y
1 − x2 − y2 
4 − x2 − y2 
9 − x2 − y2 

20. x
= y y
= −x + ye−y

2




17. x
= y y
= x2 + esin
x

7. x = y
1 − x − y  y = −x
1 − x − y  2

2

2

In each of Problems 10 through 13, use the PoincaréBendixson theorem to establish the existence of a closed trajectory of the system. In each problem, find an annular region R (region between two concentric circles) about the origin such that solutions within R remain within R. To do this, check the sign of xx
+ yy
on circles bounding the annulus. Generate a phase portrait of the system and attempt to identify closed trajectories.

471

21. x
= y y
= −x3 22. x
= 9x − 5y + x
x2 + 9y2  y
= 5x + 9y − y
x2 + 9y2 

  10. x
= x − y − x x2 + y2  y
= x + y − y x2 + y2

A differential equation x
t has a periodic solution if there is a solution x = 
t and a positive number T such that 
t + T  = 
t for all t. In each of Problems 23 through 27, prove that the differential equation has a periodic solution by converting it to a system and using theorems from this section. Generate a phase portrait for this system and attempt to identify a closed trajectory, which represents a periodic solution.

11. x
= 4x − 4y − x
x2 + 9y2  y
= 4x + 4y − y
x2 + 9y2 

23. x

+
x2 − 1x
+ 2 sin
x = 0

12. x
= y x
= −x + y − y
x2 + 2y2 

24. x

+
5x4 + 9x2 − 4xx
+ sinh
x = 0

13. x
= 4x − 2y − y
4x2 + y2  y
= 2x + 4y − y
4x2 + y2 

25. x

+ x3 = 0

In each of Problems 14 through 22, determine whether the system has a closed trajectory. Generate a phase portrait for the system and attempt to find closed trajectories. 14. x
= 3x + 4xy + xy2  y
= −2y2 + x4 y 15. x
= −y + x + x
x2 + y2  y
= x + y + y
x2 + y2  16. x
= −y2  y
= 3x + 2x3

26. x

+ 4x = 0 x 27. x

+ =0 1 + x2 28. Use Bendixson’s theorem to show that the van der Pol equation does not have a closed trajectory whose graph is completely contained in any of the following regions: (a) the infinite strip −1 < x < 1, (b) the half-plane x ≥ 1, or (c) the half-plane x ≤ −1.

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PA RT

4

CHAPTER 12 Vector Differential Calculus

CHAPTER 13

Vector Analysis

Vector Integral Calculus

The next two chapters combine vector algebra and geometry with the processes of calculus to develop vector calculus, or vector analysis. We begin with vector differential calculus, and follow in the next chapter with vector integral calculus. Much of science and engineering deals with the analysis of forces—the force of water on a dam, air turbulence on a wing, tension on bridge supports, wind and weight stresses on buildings, and so on. These forces do not occur in a static state, but vary with position, time and usually a variety of conditions. This leads to the use of vectors that are functions of one or more variables. Our treatment of vectors is in two parts—vector differential calculus (Chapter 12), and vector integral calculus (Chapter 13). Vector differential calculus extends out ability to analyze motion problems from the real line to curves and surfaces in 3-space. Tools such as the directional derivative, divergence and curl of a vector, and gradient play significant roles in many applications. Vector integral calculus generalizes integration to curves and surfaces in 3-space. This will pay many dividends, including the computation of quantities such as mass, center of mass, work, and flux of a vector field, as well as physical interpretations of vector operations. The main results are the integral theorems of Green, Gauss, and Stokes, which have broad applications in such areas as potential theory and the derivation and solution of partial differential equations modeling physical processes. 473

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CHAPTER

12

VECTOR FUNCTIONS OF ONE VARIABLE VELOCITY, ACCELERATION, CURVATURE, AND TORSION VECTOR FIELDS AND STREAMLINES THE GRADIENT FIELD AND DIRECTIONAL DERIVATIVES DIVERGENCE AND

Vector Differential Calculus

12.1

Vector Functions of One Variable In vector analysis we deal with functions involving vectors. We will begin with one such class of functions.

DEFINITION 12.1

Vector Function of One Variable

A vector function of one variable is a vector, each component of which is a function of the same single variable.

Such a function typically has the appearance F
t = x
ti + y
tj + z
tk in which x
t, y
t and z
t are the component functions of F. For each t such that the components are defined, F
t is a vector. For example, if F
t = cos
ti + 2t2 j + 3tk

(12.1)

then F
0 = i, F
 = −i + 2 2 j + 3k, and F
−3 = cos
−3i + 18j − 9k. A vector function is continuous at t0 if each component function is continuous at t0 . A vector function is continuous if each component function is continuous (for those values of t for which they are all defined). For example, G
t =

1 i + ln
tk t−1

is continuous for all t > 0 with t = 1. 475

476

CHAPTER 12

Vector Differential Calculus

The derivative of a vector function is the vector function formed by differentiating each component. With the function F of (12.1), F

t = − sin
ti + 4tj + 3k A vector function is differentiable if it has a derivative for all t for which it is defined. The vector function G defined above is differentiable for t positive and different from 1, and G

t =

−1 1 i + k 2
t − 1 t

We may think of F
t as an arrow extending from the origin to
x
t y
t z
t. Since F
t generally varies with t, we must think of this arrow as having adjustable length and pivoting at the origin to swing about as
x
t y
t z
t moves. In this way the arrow sweeps out a curve in 3-space as t varies. This curve has parametric equations x = x
t, y = y
t, z = z
t. F
t is called a position vector for this curve. Figure 12.1 shows a typical such curve. z F(t) y x FIGURE 12.1

Position vector for a curve.

The derivative of the position vector is the tangent vector to this curve. To see why this is true, observe from Figure 12.2 and the parallelogram law that the vector F
t0 + t − F
t0 

z F(t0  t)
F(t0 ) F(t0 ) F(t0  t) y x FIGURE 12.2

is represented by the arrow from
x
t0  y
t0  z
t0  to
x
t0 + t y
t0 + t z
t0 + t. Since t is a nonzero scalar, the vector 1

F
t0 + t − F
t0  t is along the line between these points. In terms of components, this vector is x
t0 + t − x
t0  y
t + t − y
t0  z
t + t − z
t0  i+ 0 j+ 0 k t t t

(12.2)

12.1 Vector Functions of One Variable

477

In the limit as t → 0, this vector approaches x

t0 i + y

t0 j + z

t0 k, which is F

t0 . In this limit, the vector (12.2) moves into a position tangent to the curve at
x
t0  y
t0  z
t0 , as suggested by Figure 12.3. This leads us to interpret F

t0  as the tangent vector to the curve at
x
t0  y
t0  z
t0 . This assumes that F

t0  is not the zero vector, which has no direction. z F
(t0 ) F(t0 )

F(t0  t) y x FIGURE 12.3 F

t0  = lim

F
t0 + t − F
t0  . t

t→0

We usually represent the tangent vector F

t0  as an arrow from the point
x
t0  y
t0  z
t0  on the curve having F
t as position vector.

EXAMPLE 12.1

Let H
t = t2 i + sin
tj − t2 k. H
t is the position vector for the curve given parametrically by x
t = t2 , y
t = sin
t, z
t = −t2 , part of whose graph is given in Figure 12.4. The tangent vector is H

t = 2ti + cos
tj − 2tk The tangent vector at the origin is H

0 = j The tangent vector at
1 sin
1 −1 is H

1 = 2i + cos
1j − 2k From calculus, we know that the length of a curve given parametrically by x = x
t, y = y
t, z = z
t for a ≤ t ≤ b, is
b length =
x

t2 +
y

t2 +
z

t2  a







in which it is assumed that x , y and z
are continuous on a b. Now  F

t =
x

t2 +
y

t2 +
z

t2

478

CHAPTER 12

Vector Differential Calculus z
1.0
0.5 0.5 1.0

20 60 80

40
40

y


60
80

x

FIGURE 12.4 Part of the graph of x = t2 , y = sin
t, z = −t2 .

is the length of the tangent vector. Thus, in terms of the position vector F
t = x
ti + y
tj + z
tk,
b length = F

t dt a

The length of a curve having a tangent at each point is the integral of the length of the tangent vector over the curve.

EXAMPLE 12.2

Consider the curve given by the parametric equations x = cos
t

y = sin
t

z=

t 3

for −4 ≤ t ≤ 4. The position vector for this curve is 1 F
t = cos
ti + sin
tj + tk 3 The graph of the curve is part of a helix wrapping around the cylinder x2 + y2 = 1, centered about the z-axis. The tangent vector at any point is 1 F

t = − sin
ti + cos
tj + k 3 Figure 12.5 shows part of the helix and tangent vectors at various points. The length of the tangent vector is 


F
t =

sin2
t + cos2
t +

1 1√ = 10 9 3

The length of this curve is length =




4

−4

F

t dt =




4 −4

8 √ 1√ 10dt = 10 3 3

12.1 Vector Functions of One Variable

479

z

y x

FIGURE 12.5

Part of a circular helix and some of its tangent vectors.

Sometimes it is convenient to write the position vector of a curve in such a way that the tangent vector at each point has length 1. Such a tangent is called a unit tangent. We will show how this can be done (at least in theory) if the coordinate functions of the curve have continuous derivatives. Let F
t = x
ti + y
tj + z
tk for a ≤ t ≤ b, and suppose x
, y
and z
are continuous. Define the real-valued function
t s
t = F

 d a

As suggested by Figure 12.6, s
t is the length of the part of the curve from its initial point
x
a y
a z
a to
x
t y
t z
t. As t moves from a to b, s
t increases from s
a = 0 to s
b = L, which is the total length of the curve. By the fundamental theorem of calculus, s
t is differentiable wherever F
t is continuous, and  ds = F

t = x

t2 + y

t2 + z

t2  dt Because s is strictly increasing as a function of t, we can, at least in theory, solve for t in terms of s, giving the inverse function t
s (see Figure 12.7). Now define G
s = F
t
s = x
t
si + y
t
sj + z
t
sk for 0 ≤ s ≤ L. Then G is a position function for the same curve as F. As t varies from a to b, F
t sweeps out the same curve as G
s, as s varies from 0 to L. However, G has the advantage that the tangent vector G
always has length 1. To see this, use the chain rule to compute d dt d F
t
s = F
t ds dt ds 1 1 F

t =
F

t = ds/dt F
t

G

s =

and this vector has length 1.

480

CHAPTER 12

Vector Differential Calculus z

C s

s(t)

(x(t), y(t), z(t))

s0  s(t 0 )

y

(x(a), y(a), z(a)) FIGURE 12.6

t0  s
1 (s0 )

x

Length function along a

t

FIGURE 12.7 A length function has an inverse.

curve.

EXAMPLE 12.3

Consider again the helix having position function 1 F
t = cos
ti + sin
tj + tk 3 for −4 ≤ t ≤ 4. We have already calculated F

t =

1√ 10. 3

Therefore the length function along this curve is
t 1√ 1√ s
t = 10d = 10
t + 4 3 −4 3 Solve for the inverse function to write 3 t = t
s = √ s − 4 10 Substitute this into the position vector to define   3 G
s = F
t
s = F √ s − 4 10       3 3 3 1 = cos √ s − 4 i + sin √ s − 4 j + √ s − 4 k 3 10 10 10       3 3 1 4 = cos √ s i + sin √ s j + √ s −  k 3 10 10 10 Now compute

    3 3 3 1 3 G
s = − √ cos √ s i + √ sin √ s j + √ k 10 10 10 10 10


This is a tangent vector to the helix, and it has length 1. Assuming that the derivatives exist, then 1. F
t + G
t
= F

t + G

t. 2. f
tF
t
= f

tF
t + f
tF

t if f is a differentiable real-valued function. 3. F
t · G
t
= F

t · G
t + F
t · G

t.

12.2 Velocity, Acceleration, Curvature and Torsion

481

4. F
t × G
t
= F

t × G
t + F
t × G

t. 5. F
f
t
= f

tF

f
t. Items (2), (3) and (4) are all “product rules”. Rule (2) is for the derivative of a product of a scalar function with a vector function; (3) is for the derivative of a dot product; and (4) is for the derivative of a cross product. In each case the rule has the same form as the familiar calculus formula for the derivative of a product of two functions. However, in (4), order is important, since F × G = −G × F. Rule (5) is a vector version of the chain rule. In the next section we will use vector functions to develop the concepts of velocity and acceleration, which we will apply to the geometry of curves in 3–space.

SECTION 12.1

PROBLEMS

In each of Problems 1 through 8, compute the requested derivative (a) by carrying out the vector operation and differentiating the resulting vector or scalar, and (b) by using one of the differentiation rules (1) through (5) stated at the end of this section. 1. F
t = i+3t j+2tk f
t = 4 cos
3t
d/dt f
tF
t 2

2. F
t = ti − 3t2 k G
t = i + cos
tk
d/dt F
t · G
t 3. F
t = ti + j + 4k G
t = i − cos
tj + tk
d/dt F
t × G
t 4. F
t = sinh
tj − tk G
t = ti + t2 j − t2 k
d/dt F
t × G
t 5. F
t = ti − cosh
tj + et k f
t = 1 − 2t3 
d/dt f
tF
t 6. F
t = ti − tj + t2 k G
t = sin
ti − 4j + t3 k
d/dt F
t · G
t

12.2

7. F
t = −9i + t2 j + t2 k G
t = et i
d/dt F
t × G
t 8. F
t = −4 cos
tk G
t = −t2 i + 4 sin
tk
d/dt F
t · G
t In each of Problems 9, 10, and 11, (a) write the position vector and tangent vector for the curve whose parametric equations are given, (b) find a length function s
t for the curve, (c) write the position vector as a function of s, and (d) verify that the resulting position vector has a derivative of length 1. 9. x = sin
t y = cos
t z = 45t
0 ≤ t ≤ 2 10. x = y = z = t3 
−1 ≤ t ≤ 1 11. x = 2t2  y = 3t2  z = 4t2 
1 ≤ t ≤ 3 12. Let F
t = x
ti + y
tj + z
tk. Suppose x, y and z are differentiable functions of t. Think of F
t as the position function of a particle moving along a curve in 3-space. Suppose F × F
= O. Prove that the particle always moves in the same direction.

Velocity, Acceleration, Curvature and Torsion Imagine a particle moving along a path having position vector F
t = x
ti + y
tj + z
tk as t varies from a to b. We want to relate F to the dynamics of the particle. For calculations we will be doing, we assume that x, y and z are twice differentiable. We will also make use of the distance function along the curve, s
t =




t a

F

 d

482

CHAPTER 12

Vector Differential Calculus

DEFINITION 12.2

Velocity, Speed

1. The velocity v
t of the particle at time t is defined to be v
t = F

t 2. The speed v
t of the particle at time t is the magnitude of the velocity.

Velocity is therefore a vector, having magnitude and direction. If v
t is not the zero vector, then the velocity is tangent to the curve of motion of the particle. Thus, at any instant the particle may be thought of as moving in the direction of the tangent to the path of motion. The speed at time t is a real-valued function, given by v
t = v
t = F

t =

ds  dt

This is consistent with the familiar idea of speed as the rate of change of distance (along the path of motion) with respect to time.

DEFINITION 12.3

Acceleration

The acceleration a
t of the particle is the rate of change of the velocity with respect to time: a
t = v

t

Alternatively, a
t = F


t As with velocity, acceleration is a vector.

EXAMPLE 12.4

Let F
t = sin
ti + 2e−t j + t2 k The path of the particle is the curve whose parametric equations are x = sin
t

y = 2e−t 

z = t2 

Part of the graph of this curve is shown in Figure 12.8. The velocity and acceleration are, respectively, v
t = cos
ti − 2e−t j + 2tk and a
t = − sin
ti + 2e−t j + 2k

12.2 Velocity, Acceleration, Curvature and Torsion

483

z 80 60

40


1

20 0 1

x

1 2 y

FIGURE 12.8 Part of the graph of x = sin
t, y = 2e−t , z = t2 .

The speed of the particle is v
t =



cos2
t + 4e−2t + 4t2 

If F

t is not the zero vector, then this vector is tangent to the curve at
x
t y
t z
t. We may obtain a unit tangent by dividing this vector by its length: T
t =

1 F

t

F

t =

1 F

t ds/dt

Equivalently, T
t =

1 1 v
t v
t = v
t v
t

the velocity divided by the speed. If arc length s along the path of motion is used as the parameter in the position function, then we have seen that F

s = 1 automatically, so in this case the speed is identically 1 and this unit tangent is just the velocity vector. We will use this unit tangent to define a function that quantifies the “amount of bending” of a curve at a point.

DEFINITION 12.4

Curvature

The curvature % of a curve is the magnitude of the rate of change of the unit tangent with respect to arc length along the curve: # # # dT # # %
s = # # ds # 

The definition is motivated by the intuition (Figure 12.9) that the more a curve bends at a point, the faster the tangent vector is changing there. If the unit tangent vector has been written using s as parameter, then computing T

s is straightforward. More often, however, the unit tangent is parametrized by some other variable, and then the derivative defining the curvature must be computed by the chain rule: # # # dT dt # # # %
t = # dt ds #

484

CHAPTER 12

Vector Differential Calculus z

y

x FIGURE 12.9 Increasing curvature corresponds to an increasing rate of change of the tangent vector along the curve.

This gives the curvature as a function of the parameter used in the position vector. Since dt 1 1 = =  ds ds/dt F

t we often write %
t =

1 F

t

T

t 

EXAMPLE 12.5

Consider a straight line having parametric equations x = a + bt

y = c + dt

z = e + ht

in which a, b, c, d, e and h are constants. The position vector of this line is F
t =
a + bti +
c + dtj +
e + htk We will compute the curvature using equation (12.3). First, F

t = bi + dj + hk so F

t =



b 2 + d 2 + h2 

The unit tangent vector is T
t =

1 F

t

=√

F

t

1 b 2 + d 2 + h2


bi + dj + hk

(12.3)

12.2 Velocity, Acceleration, Curvature and Torsion

485

This is a constant vector, so T

t = O and the curvature is %
t =

1 T

t = 0 F

t

This is consistent with our intuition that a straight line should have curvature zero.

EXAMPLE 12.6

Let C be the circle of radius 4 about the origin in the plane y = 3. Using polar coordinates, this curve has parametric equations x = 4 cos


y = 3

z = 4 sin


for 0 ≤ t ≤ 2. The circle has position vector F
 = 4 cos
i + 3j + 4 sin
k Then F

 = −4 sin
i + 4 cos
k so F

 = 4 The unit tangent is T
 =

1

−4 sin
i + 4 cos
k = − sin
i + cos
k 4

Then T

 = − cos
i − sin
k The curvature is %
 =

1 1 T

 =  4 F



The curvature of this circle is constant, again consistent with intuition. One can show that a circle of radius r has curvature 1/r. Not only does a circle have constant curvature, but this curvature also decreases the larger the radius is chosen, as we should expect.

EXAMPLE 12.7

Let C have parametric representation x = cos
t + t sin
t

y = sin
t − t cos
t

z = t2

for t > 0. Figure 12.10 shows part of the graph of C. We will compute the curvature. We can write the position vector F
t = cos
t + t sin
t i + sin
t − t cos
t j + t2 k

486

CHAPTER 12

Vector Differential Calculus z

–8

–6

–4

8

80 60 40 20 –2

0 –2

4

4

6

–4

6

8

x

y

FIGURE 12.10 Part of the graph of x = cos
t + t sin
t, y = sin
t− t cos
t, z = t2 .

A tangent vector is given by F

t = t cos
ti + t sin
tj + 2tk and F

t =

√

5t

Next, the unit tangent vector is T
t =

1 F

t

1 F

t = √ cos
ti + sin
tj + 2k  5

Compute 1 T

t = √ − sin
ti + cos
tj  5 We can now use equation (12.3) to compute %
t =

1 F

t

1 =√ 5t



T

t

 1 2 1 sin
t + cos2
t = 5 5t

for t > 0. We will now introduce another vector of interest in studying motion along a curve.

DEFINITION 12.5

Unit Normal Vector

Using arc length s as parameter on the curve, the unit normal vector N
s is defined by N
s = provided that %
s = 0.

1
T
s %
s

12.2 Velocity, Acceleration, Curvature and Torsion

487

The name given to this vector is motivated by two observations. First, N
s is a unit vector. Since %
s = T

s, then N
s =

1 T

s = 1 T

s

Second, N
s is orthogonal to the unit tangent vector. To see this, begin with the fact that T
s is a unit tangent, hence T
s = 1. Then T
s2 = T
s · T
s = 1 Differentiate this equation to get T

s · T
s + T
s · T

s = 2T
s · T

s = 0 hence T
s · T

s = 0 which means that T
s is orthogonal to T

s. But N
s is a positive scalar multiple of T

s, and so is in the same direction as T

s. We conclude that N
s is orthogonal to T
s At any point of a curve with differentiable coordinate functions (not all vanishing for the same parameter value), we may now place a tangent vector to the curve, and a normal vector that is perpendicular to the tangent vector (Figure 12.11). z

T N

y x FIGURE 12.11 Tangent and normal vector to a curve at a point.

EXAMPLE 12.8

Consider again the curve with position function F
t = cos
t + t sin
t i + sin
t − t cos
t j + t2 k for t > 0. In Example 12.7 we computed the unit tangent and the curvature as functions of t. We will write the position vector as a function of arc length, and compute the unit tangent T
s and the unit normal N
s.

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First, using F

t =

√

5t from Example 12.7,

s
t =




t 0




F
 d =

Solve for t as a function of s:

in which  =


t√ 0

√ 5d =

5 2 t  2

√ √ 2√ t = 1/4 s =  s 5

√ . 1/4 2 5 . In terms of s, the position vector is

G
s = F
t
s

√ √ √ 

√ √ √  = cos
 s +  s sin
 s i + sin
 s −  s cos
 s j + 2 sk The unit tangent is T
s = G

s √ √ 1 1 = 2 cos
 si + 2 sin
 sj + 2 k 2 2 This vector does indeed have length 1: 1 5 T
s = 4 + 4 = 4 4 2

 √ 4 2 = 1 51/4

Now √ √ 3 3 T

s = − √ sin
 si + √ cos
 sj 4 s 4 s so the curvature is  %
s = T

s =

6 16s

1/2

1 = √ 4 s



21/2 51/4

3 =

1 53/4

1 1 √ √ 2 s

√ for s > 0. Since s = 5t2 /2, then in terms of t we have % = 1/5t, consistent with Example 12.7. Now compute the unit normal vector √

√ √ 4 s 3 1
3 N
s = T
s = 3 − √ sin
 si + √ cos
 sj %
s  4 s 4 s √ √ = − sin
 si + cos
 sj This is a unit vector orthogonal to T
s.

12.2.1 Tangential and Normal Components of Acceleration At any point on the trajectory of a particle, the tangent and normal vectors are orthogonal. We will now show how to write the acceleration at a point as a linear combination of the tangent and normal vectors there: a = aT T + aN N This is illustrated in Figure 12.12.

12.2 Velocity, Acceleration, Curvature and Torsion

489

z a a TT aNN y a  a TT  a N N

x

FIGURE 12.12

THEOREM 12.1

a=

dv T + v2 %N dt

Thus aT = dv/dt and aN = v2 %. The tangential component of the acceleration is the derivative of the speed, while the normal component is the curvature at the point, times the square of the speed there. Proof

First observe that T
t =

1 1 F

t = v v F

t

Therefore v = vT Then d dv v= T + vT

t dt dt ds dT dv T+v = dt dt ds dv T + v2 T

s = dt dv T + v2 %N = dt

a=

Here is one use of this decomposition of the acceleration. Since T and N are orthogonal unit vectors, then a2 = a · a =
aT T + aN N ·
aT T + aN N = a2T T · T + 2aT aN T · N + a2N N · N = a2T + a2N  From this, whenever two of a, aT and aN are known, we can compute the third quantity.

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EXAMPLE 12.9

Return again to the curve C having position function F
t = cos
t + t sin
t i + sin
t − t cos
t j + t2 k for t > 0. We will compute the tangential and normal components of the acceleration. First, v
t = F

t = t cos
ti + t sin
tj + 2tk so the speed is v
t = F

t =

√

5t

The tangential component of the acceleration is therefore dv √ aT = = 5 dt a constant for this curve. The acceleration vector is a = v
= cos
t − t sin
ti + sin
t + t cos
tj + 2k and a routine calculation gives a =



5 + t2 

Then a2N = a2 − a2T = 5 + t2 − 5 = t2  Since t > 0, the normal component of acceleration is aN = t. The acceleration may therefore be written as √ a = 5T + tN If we know the normal component aN and the speed v, it is easy to compute the curvature, since aN = t = %v2 = 5t2 % implies that 1  5t as we computed in Example 12.7 directly from the tangent vector. Now the unit tangent and normal vectors are easy to compute in terms of t. First, %=

1 1 T
t = v = √ cos
ti + sin
tj + 2k v 5 This is usually easy to compute, since v = F

t is a straightforward calculation. But in addition, we now have the unit normal vector (as a function of t) 1
1 dt dT 1
T
s = = T
t % % ds dt %v 5t 1 = √ √ − sin
ti + cos
tj = − sin
ti + cos
tj 5t 5

N
t =

This calculation does not require the explicit computation of s
t and its inverse function.

12.2 Velocity, Acceleration, Curvature and Torsion

12.2.2

491

Curvature as a Function of t

Equation (12.3) gives the curvature in terms of the parameter t used for the position function: %
t =

1 T

t  F

t

This is a handy formula because it does not require introduction of the distance function s
t. We will now derive another expression for the curvature that is sometimes useful in calculating % directly from the position function.

THEOREM 12.2

Let F be the position function of a curve, and suppose that the components of F are twice differentiable functions. Then %=

F
× F

 F
3/2



This states that the curvature is the magnitude of the cross product of the first and second derivatives of F, divided by the cube of the length of F
. Proof

First write 

ds a = aT T + % dt

2 N

Take the cross product of this equation with the unit tangent vector:  2  2 ds ds T × a = aT T × T + %
T × N = %
T × N dt dt since the cross product of any vector with itself is the zero vector. Then  2 ds T × a = % T × N dt  2 ds =% T N sin
 dt where  is the angle between T and N. But these are orthogonal unit vectors, so  = /2 and T = N = 1. Therefore  2 ds  T × a = % dt Then %=

T × a 
ds/dt2

But T = F
/ F
, a = T
= F

, and ds/dt = F
, so # #


# 1 # # 1 F
× F

# = F × F   %= 2 3 # #
F
 F  F


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EXAMPLE 12.10

Let C have position function F
t = t2 i − t3 j + tk We want the curvature of C. Compute F

t = 2ti − 3t2 j + k F


t = 2i − 6tj and   i 


F × F =  2t  2

j −3t2 −6t

k 1 0

     

= 6ti + 2j − 6t2 k The curvature is %
t = =

# # #6ti + 2j − 6t2 k# 2ti − 3t2 j + k3 √ 36t2 + 4 + 36t4
4t2 + 9t4 + 13/2



12.2.3 The Frenet Formulas If T and N are the unit tangent and normal vectors, the vector B = T × N is also a unit vector, and is orthogonal to T and N. At any point on the curve where these three vectors are defined and nonzero, the triple T, N, B forms a right-handed rectangular coordinate system (as in Figure 12.13). We can in effect put an x y z coordinate system at any point P on C, with the positive x-axis along T, the positive y-axis along N, and the positive z axis along B. Of course, this system twists and changes orientation in space as P moves along C (Figure 12.14). Since N =
1/%T

s, then dT = %N ds

z

N

B

T

B

T y

x FIGURE 12.13

N

N B FIGURE 12.14

N B T

T

C

12.3 Vector Fields and Streamlines

493

Further, it can be shown that there is a scalar-valued function  such that dN = −%T + B ds and dB = −N ds These three equations are called the Frenet formulas. The scalar quantity 
s is the torsion of C at
x
s y
s z
s. If we look along C at the coordinate system formed at each point by T, N and B, the torsion measures how this system twists about the curve as the point moves along the curve.

SECTION 12.2

PROBLEMS

In each of Problems 1 through 10, a position vector is given. Determine the velocity, speed, acceleration, tangential and normal components of the acceleration, the curvature, the unit tangent, unit normal and binormal vectors. 1. F = 3ti − 2j + t2 k

8. F = ln
t
i − j + 2k 9. F = t2 i + t2 j + !t2 k 10. F = 3t cos
tj − 3t sin
tk 11. Suppose we are given the position vector of a curve and find that the unit tangent vector is a constant vector. Prove that the curve is a straight line.

2. F = t sin
ti + t cos
tj + k 3. F = 2ti − 2tj + tk

12. It is easy to verify that the curvature of any straight line is zero. Suppose C is a curve with twice differentiable coordinate functions, and curvature zero. Does it follow that C is a straight line?

4. F = et sin
ti − j + et cos
tk 5. F = 3e−t
i + j − 2k 6. F =  cos
ti + tj +  sin
tk

12.3

7. F = 2 sinh
tj − 2 cosh
tk

Vector Fields and Streamlines We now turn to the analysis of vector functions of more than one variable.

DEFINITION 12.6

Vector Field

A vector field in 3-space is a 3-vector whose components are functions of three variables. A vector field in the plane is a 2-vector whose components are functions of two variables.

A vector field in 3-space has the appearance G
x y z = f
x y zi + g
x y zj + h
x y zk and, in the plane, K
x y = f
x yi + g
x yj

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The term “vector field” is geometrically motivated. At each point P for which a vector field G is defined, we can represent the vector G
P as an arrow from P. If is often useful in working with a vector field G to draw arrows G
P at points through the region where G is defined. This drawing is also referred to as a vector field (think of arrows growing at points). The variations in length and orientation of these arrows gives some sense of the flow of the vector field, and its variations in strength, just as a direction field helps us visualize trajectories of a system X
= F
X of differential equations. Figures 12.15 and 12.16 show the vector fields G
x y = xyi +
x − yj and H
x y = y cos
xi +
x2 − y2 j, respectively, in the plane. Figures 12.17, 12.18 and 12.19 show the vector fields F
x y z = cos
x + yi − xj +
x − zk, Q
x y z = −yi + zj +
x + y + zk and M
x y z = cos
xi + e−x sin
yj +
z − yk, respectively, in 3-space.

y

y 1

1 0
1
1

1

x

FIGURE 12.15 Representation of the vector field G
x y = xyi +
x − yj.


1 0
1

x

1

FIGURE 12.16 Representation of the vector field H
x y = y cos
xi+
x2 − y2 j.

y
1 1 z
1

y 1

x

1

z

FIGURE 12.17 Representation of the vector field F
x y z = cos
x + yi − xj +
x − zk as arrows in planes z = constant.


1
1

1

x

FIGURE 12.18 Representation of the vector field Q
x y z = −yi + zj +
x + y + zk in planes z = constant.

12.3 Vector Fields and Streamlines

495

z

x 1 y

0

1


1
1

FIGURE 12.19 The vector field M
x y z = cos
xi + e−x sin
yj +
z − yk in planes z = constant.

A vector field is continuous if each of its component functions is continuous. A partial derivative of a vector field is the vector field obtained by taking the partial derivative of each component function. For example, if F
x y z = cos
x + yi − xj +
x − zk then

F = Fx = − sin
x + yi − j + k

x

F = Fy = − sin
x + yi

y

and

F = Fz = −k

z

If G
x y = xyi +
x − yj then

G = Gx = yi + j

x

and

G = Gy = xi − j

y

Streamlines of a vector field F are curves with the property that, at each point
x y z, the vector F
x y z is tangent to the curve through this point.

DEFINITION 12.7

Streamlines

Let F be a vector field, defined for all
x y z in some region  of 3-space. Let
be a set of curves with the property that, through each point P of , there passes exactly one curve from
. The curves in
are streamlines of F if, at each
x y z in , the vector F
x y z is tangent to the curve in
passing through
x y z.

Streamlines are also called flow lines or lines of force, depending on context. If F is the velocity field for a fluid, the streamlines are often called flow lines (paths of particles in the

496

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fluid). If F is a magnetic field the streamlines are called lines of force. If you put iron filings on a piece of cardboard and then hold a magnet underneath, the filings will be aligned by the magnet along the lines of force of the field. Given a vector field F, we would like to find the streamlines. This is the problem of constructing a curve through each point of a region, given the tangent to the curve at each point. Figure 12.20 shows typical streamlines of a vector field, together with some of the tangent vectors. We want to determine the curves from the tangents. z

y x FIGURE 12.20

Streamlines of a vector field.

To solve this problem, suppose C is a streamline of F = f i + gj + hj. Let C have parametric equations x = x
 y = y
 z = z
 The position vector for this curve is R
 = x
i + y
j + z
k Now R

 = x

i + y

j + z

k is tangent to C at
x
 y
 z
. But for C to be a streamline of F, F
x
 y
 z
 is also tangent to C at this point, hence must be parallel to R

. These vectors must therefore be scalar multiples of each other. For some scalar t (which may depend on ), R

 = tF
x
 y
 z
 But then dx dy dz i + j + k = tf
x
 y
 z
i + tg
x
 y
 z
j + th
x
 y
 z
k d d d This implies that dx dy dz = tf = tg = th d d d

(12.4)

Since f , g and h are given functions, these equations constitute a system of differential equations for the coordinate functions of the streamlines. If f , g and h are nonzero, then t can be eliminated to write the system in differential form as dx dy dz = =  f g h

(12.5)

12.3 Vector Fields and Streamlines

497

EXAMPLE 12.11

We will find the streamlines of the vector field F = x2 i + 2yj − k. The system (12.4) is dy dx dz = tx2  = 2ty = −t d d d If x and y are not zero, this can be written in the form of equations (12.5): dx dy dz =  = x2 2y −1 These equations can be solved in pairs. To begin, integrate dx = −dz x2 to get 1 − = −z + c x in which c is an arbitrary constant. Next, integrate dy = −dz 2y to get 1 ln y = −z + k 2 It is convenient to express two of the variables in terms of the third. If we solve for x and y in terms of z, we get x=

1  z−c

y = ae−2z 

in which a = e2k is an arbitrary (positive) constant. This gives us parametric equations of the streamlines, with z as parameter: x=

1  z−c

y = ae−2z 

z = z

To find the streamline through a particular point, we must choose c and a appropriately. For example, suppose we want the streamline through
−1 6 2. Then z = 2 and we need to choose c and a so that −1 =

1 and 6 = ae−4  2−c

Then c = 3 and a = 6e4 , so the streamline through
−1 6 2 has parametric equations x=

1  z−3

y = 6e4−2z 

A graph of this streamline is shown in Figure 12.21.

z = z

498

CHAPTER 12

Vector Differential Calculus z 4 2 4000 x 10


2

8000

12000 16000 y

FIGURE 12.21 Part of the graph of x = 1/
z − 3, y = 6e4−2z , z = z.

EXAMPLE 12.12

Suppose we want the streamlines of F
x y z = −yi + zk. Here the i component is zero, so we must begin with equations (12.4), not (12.5). We have dx = 0 d

dy = −ty d

dz = tz d

The first equation implies that x = constant. This simply means that all the streamlines are in planes parallel to the y z plane. The other two equations yield dy dz = −y z and an integration gives − ln
y + c = ln
z Then ln
zy = c implying that zy = k in which k is constant. The streamlines are given by the equations x = c

k z=  y

in which c and k are arbitrary constants and y is the parameter. For example, to find the streamline through
−4 1 7, choose c = −4 and k so that z = k/y passes through y = 1 z = 7. We need 7=

k = k 1

and the streamline has equations 7 x = −4 z =  y The streamline is a hyperbola in the plane x = −4 (Figure 12.22).

12.4 The Gradient Field and Directional Derivatives


8


4.0

x

499

z
4 20


4.2
4.4
20

0 4 8

y


3.8
3.6

FIGURE 12.22 Part of the graph of x = −4,

z = 7/y.

SECTION 12.3

PROBLEMS

In each of Problems 1 through 5, compute the two first partial derivatives of the vector field and make a diagram in which each indicated vector is drawn as an arrow from the point at which the vector is evaluated. 1. G
x y = 3xi − 4xyj G
0 1 G
1 3 G
1 4 G
−1 −2 G
−3 2 2. G
x y = ex i − 2x2 yj G
0 0 G
0 1 G
2 −3 G
−1 −3 3. G
x y = 2xyi+cos
xj G
/2 0 G
0 0 G
−1 1 G
 −3 G
−/4 −2 4. G
x y = sin
2xyi +
x2 + yj G
−/2 0 G
0 2 G
/4 4 G
1 1 G
−2 1 5. G
x y = 3x2 i+
x−2yj G
1 −1 G
0 2 G
−3 2 G
−2 −2 G
2 5 In each of Problems 6 through 10, compute the three first partial derivatives of the vector field. 6. F = exy i − 2x2 yj + cosh
z + yk 7. F = 4z2 cos
xi − x3 yzj + x3 yk

12.4

8. F = 3xy3 i + ln
x + y + zj + cosh
xyzk 9. F = −z4 sin
xyi + 3xy4 zj + cosh
z − xk 10. F =
14x − 2yi +
x2 − y2 − z2 j + 5xyk In each of Problems 11 through 16, find the streamlines of the vector field, then find the particular streamline through the given point. 11. F = i − y2 j + zk
2 1 1 12. F = i − 2j + k
0 1 1 13. F =
1/xi + ex j − k
2 0 4 14. F = cos
yi + sin
xj
/2 0 −4 15. F = 2ez j − cos
yk
3 /4 0 16. F = 3x2 i − yj + z3 k
2 1 6 17. Construct a vector field whose streamlines are straight lines. 18. Construct a vector field in the x y plane whose streamlines are circles about the origin.

The Gradient Field and Directional Derivatives Let 
x y z be a real-valued function of three variables. In the context of vectors, such a function is called a scalar field. We will define an important vector field manufactured from .

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DEFINITION 12.8

Gradient

The gradient of a scalar field  is the vector field & given by & =





 i + j + k

x

y

z

wherever these partial derivatives are defined.

The symbol & is read “del phi”, and & is called the del operator. It operates on a scalar field to produce a vector field. For example, if 
x y z = x2 y cos
yz, then & = 2xy cos
yzi + x2 cos
yz − x2 yz sin
yzj − x2 y2 sin
yzk The gradient field evaluated at a point P is denoted &
P. For the gradient just computed, &
1 −1 3 = −2 cos
3i + cos
3 − 3 sin
3j + sin
3k If  is a function of just x and y, then & is a vector in the x y plane. For example, if 
x y =
x − y cos
y, then &
x y = cos
yi + − cos
y −
x − y sin
yj At
2  this gradient is &
2  = −i + j The gradient has the obvious properties &
 +  = & + & and, if c is a number, then &
c = c& We will now define the directional derivative, and relate this to the gradient. Suppose 
x y z is a scalar field. Let u = ai + bj + ck be a unit vector (length 1). Let P0 =
x0  y0  z0 . Represent u as an arrow from P0 , as in Figure 12.23. We want to define a quantity that measures the rate of change of 
x y z as
x y z varies from P0 , in the direction of u. To do this, notice that, if t > 0, then the point P 
x0 + at y0 + bt z0 + ct is on the line through P0 in the direction of u. Further, the distance from P0 to P along this direction is exactly t, because the vector from P0 to P is
x0 + at − x0 i +
y0 + bt − y0 j +
z0 + ct − z0 k

(x, y, z)  (x 0  at, y0  bt, z 0  ct) t0 P0 (x 0 , y 0 , z 0 )

u

FIGURE 12.23

12.4 The Gradient Field and Directional Derivatives

501

and this is just tu. The derivative d 
x + at y + bt z + ct dt is the rate of change of 
x + at y + bt z + ct with respect to this distance t, and

d 
x + at y + bt z + ct dt t=0 is this rate of change evaluated at P0 . This derivative gives the rate of change of 
x y z at P0 in the direction of u. We will summarize this discussion in the following definition. DEFINITION 12.9

Directional Derivative

The directional derivative of a scalar field  at P0 in the direction of the unit vector u is denoted Du 
P0 , and is given by

d 
x + at y + bt z + ct Du 
P0  =  dt t=0 We usually compute a directional derivative using the following. THEOREM 12.3

If  is a differentiable function of two or three variables, and u is a constant unit vector, then Du 
P0  = &
P0  · u Proof

Let u = ai + bj + ck. By the chain rule, d





 
x + at y + bt z + ct = a + b + c dt

x

y

z

Since
x0 + at y0 + bt z0 + ct =
x0  y0  z0  when t = 0, then Du 
P0  =






P a +
P0 b +
P0 c

x 0

y

z

= &
P0  · u

EXAMPLE 12.13

Let 
x y z = x2 y − xez , P0 =
2 −1  and u = 
x y z at P0 in the direction of u is Du 
2 −1  = &
2 −1  · u

√1 6


i − 2j + k. Then the rate of change of

  1 2 1 = x
2 −1  √ + y
2 −1  − √ + z
2 −1  √ 6 6 6 

1  = √ 2xy − ez 
2−1 − 2 x2
−21 + −xez 
2−1 6 1 −3 = √
−4 − e − 8 − 2e  = √
4 + e   6 6

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CHAPTER 12

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In working with directional derivatives, care must be taken that the direction is given by a unit vector. If a vector w of length other than 1 is used to specify the direction, then use the unit vector w/ w in computing the directional derivative. Of course, w and w/ w have the same direction. A unit vector is used with directional derivatives so that the vector specifies only direction, without contributing a factor of magnitude. Suppose now that 
x y z is defined at least for all points within some sphere about P0 . Imagine standing at P0 and looking in various directions. We may see 
x y z increasing in some, decreasing in others, perhaps remaining constant in some directions. In what direction does 
x y z have its greatest, or least, rate of increase from P0 ? We will now show that the gradient vector &
P0  points in the direction of maximum rate of increase at P0 , and −&
P0  in the direction of minimum rate of increase. THEOREM 12.4

Let  and its first partial derivatives be continuous in some sphere about P0 , and suppose that &
P0  = O. Then 1. At P0 , 
x y z has its maximum rate of change in the direction of &
P0 . This maximum rate of change is &
P0 . 2. At P0 , 
x y z has its minimum rate of change in the direction of −&
P0 . This minimum rate of change is − &
P0 . Proof

Let u be any unit vector. Then Du 
P0  = &
P0  · u = &
P0  u cos
 = &
P0  cos


because u has length 1.  is the angle between u and &
P0 . The direction u in which  has its greatest rate of increase from P0 is the direction in which this directional derivative is a maximum. Clearly the maximum occurs when cos
 = 1, hence when  = 0. But this occurs when u is along &
P0 . Therefore this gradient is the direction of maximum rate of change of 
x y z at P0 . This maximum rate of change is &
P0 . For (2), observe that the directional derivative is a minimum when cos
 = −1, hence when  = . This occurs when u is opposite &
P0 , and this minimum rate of change is − &
P0 .

EXAMPLE 12.14

Let 
x y z = 2xz + ey z2 . We will find the maximum and minimum rates of change of 
x y z from
2 1 1. First, &
x y z = 2zi + ey z2 j +
2x + 2zey k so &
P0  = 2i + ej +
4 + 2ek The maximum rate of increase of 
x y z at
2 1 1 is in the direction of this gradient, and this maximum rate of change is  4 + e2 +
4 + 2e2  The  minimum rate of increase is in the direction of −2i − ej −
4 + 2ek, and is − 4 + e2 +
4 + 2e2 .

12.4 The Gradient Field and Directional Derivatives

12.4.1

503

Level Surfaces, Tangent Planes and Normal Lines

Depending on the function  and the constant k, the locus of points 
x y z = k may form a surface in 3-space. Any such surface is called a level surface of . For example, if 
x√ y z = x2 + y2 + z2 and k > 0, then the level surface 
x y z = k is a sphere of radius k. If k = 0 this locus is just a single point, the origin. If k < 0 this locus is empty. There are no points whose coordinates satisfy this equation. The level surface 
x y z = 0 of 
x y z = z − sin
xy is shown from three perspectives in Figures 12.24 (a), (b) and (c).

y

2

z 1

-2

x 1

0

-1

2

-1 -2

(a)

y 2 1

z

-2 -1

0 1 -1

2 x

-2

(b) FIGURE 12.24 Different perspectives of graphs of

z = sin
xy.

504

CHAPTER 12

Vector Differential Calculus

z -2 -1 -1 0

-2

1 2 y

1

2

x

(c) FIGURE 12.24

(Continued).

Now consider a point P0 
x0  y0  z0  on a level surface 
x y z = k. Assume that there are smooth (having continuous tangents) curves on the surface passing through P0 , such as C1 and C2 in Figure 12.25. Each such curve has a tangent vector at P0 . These tangent vectors determine a plane at P0 , called the tangent plane to the surface at P0 . A vector normal (perpendicular) to at P0 , in the sense of being normal to each of these tangent vectors, is called a normal vector, or just normal, to the surface at P0 . We would like to be able to determine the tangent plane and normal to a surface at a point. Recall that we can find the equation of a plane through a given point if we are given a normal vector to the plane. Thus the normal vector is the key to finding the tangent plane. z

∏

P0

C1 C2

y  (x, y, z) = k x FIGURE 12.25 Tangents to curves on the surface through P0 determine the tangent plane .

THEOREM 12.5

Gradient As a Normal Vector

Let  and its first partial derivatives be continuous. Then &
P is normal to the level surface 
x y z = k at any point P on this surface at which this gradient vector is nonzero.

12.4 The Gradient Field and Directional Derivatives

505

We will outline an argument suggesting why this is true. Consider a point P0 
x0  y0  z0  on the level surface 
x y z = k. Suppose a smooth curve C on this surface passes through P0 , as in Figure 12.26 (a). Suppose C has parametric equations x = x
t y = y
t z = z
t

z  (P0 )  (P0 ) P0 C y (x, y, z)  constant

x

Tangent plane to the level surface at P0

P0

Level surface (x, y, z)  k

(a)

(b)

FIGURE 12.26 
P0  is normal to the level surface at P0 .

Since P0 is on this curve, for some t0 , x
t0  = x0  y
t0  = y0  z
t0  = z0  Further, since the curve lies on the level surface, then 
x
t y
t z
t = k for all t. Then d 
x
t y
t z
t = 0 dt = x x

t + y y

t + z z

t = & · x

ti + y

tj + z

tk Now x

ti + y

tj + z

tk = T
t is a tangent vector to C. In particular, letting t = t0 , then T
t0  is a tangent vector to C at P0 , and we have &
P0  · T
t0  = 0 This means that &
P0  is normal to the tangent to C at P0 . But this is true for any smooth curve on the surface and passing through P0 . Therefore &
P0  is normal to the surface at P0 (Figure 12.26 (b)). Once we have this normal vector, finding the equation of the tangent plane is straightforward. If
x y z is any other point on the tangent plane (Figure 12.27), then the vector
x − x0 i +
y − y0 j +
z − z0 k is in this plane, hence is orthogonal to the normal vector. Then &
P0  ·
x − x0 i +
y − y0 j +
z − z0 k = 0 Then






P 
x − x0  +
P0 
y − y0  +
P0 
z − z0  = 0

x 0

y

z

(12.6)

506

CHAPTER 12

Vector Differential Calculus

(x, y, z)

z

P0 : (x 0 , y 0 , z 0 )

 (P0 )

(x, y, z)  constant y x FIGURE 12.27


P0 ·
x − x0 i +
y − y0 j +
z − z0 k = 0.

This equation is satisfied by every point on the tangent plane. Conversely, if
x y z satisfies this equation, then
x − x0 i +
y − y0 j +
z − z0 k is normal to the normal vector, hence lies in the tangent plane, implying that
x y z is a point in this plane. We call equation (12.6) the equation of the tangent plane to 
x y z = k at P0 .

EXAMPLE 12.15

 Consider the level surface 
x y z = z − x2 + y2 = 0. This surface is the cone shown √ in Figure 12.28. We will find the normal vector and tangent plane to this surface at
1 1 2. First compute the gradient vector: & = − 

x x2 + y 2

i− 

y

j + k

x2 + y 2

provided that x and y are not both zero. Figure 12.29 shows & at a point on the cone determined by the position vector  R
x y z = xi + yj + x2 + y2 k Then √ 1 1 &
1 1 2 = − √ i − √ j + k 2 2

z z 

z  x 2  y 2

R

y x FIGURE 12.28

P0 y

x FIGURE  12.29 Cone z = x2 + y2 and normal at P0 .

12.4 The Gradient Field and Directional Derivatives

507

√ This is the normal vector to the cone at
1 1 2. The tangent plane at this point has equation √ 1 1 − √
x − 1 − √
y − 1 + z − 2 = 0 2 2 or √ x + y − 2z = 0 The cone has no tangent plane or normal vector at the origin, where the surface has a “sharp point”. This is analogous to a graph in the plane having no tangent vector where it has a sharp point (for example, y = x at the origin).

EXAMPLE 12.16

Consider the surface z = sin
xy. If we let 
x y z = sin
xy − z, then this surface is the level surface 
x y z = 0. The gradient vector is & = y cos
xyi + x cos
xyj − k This vector field is shown in Figure 12.30, with the gradient vectors drawn as arrows from selected points on the surface. The tangent plane at any point
x0  y0  z0  on this surface has equation y0 cos
x0 y0 
x − x0  + x0 cos
x0 y0 
y − y0  −
z − z0  = 0

z
2


2
1

1


1

1

2 2 x y

FIGURE 12.30 Gradient field  = y cos
xyi + x cos
xyj − k represented as a vector field

on the surface z = sin
xy.

508

CHAPTER 12

Vector Differential Calculus

For example, the tangent plane at
2 1 sin
2 has equation cos
2
x − 2 + 2 cos
2
y − 1 − z + sin
2 = 0 or cos
2x + 2 cos
2y − z = 4 cos
2 − sin
2 A patch of this tangent plane is shown in Figure 12.31. Similarly, the tangent plane at
−1 −2 sin
2 has equation 2 cos
2x + cos
2y + z = −4 cos
2 + sin
2 Part of this tangent plane is shown in Figure 12.32. A straight line through P0 and parallel to &
P0  is called the normal line to the level surface 
x y z = k at P0 , assuming that this gradient vector is not zero. This idea is illustrated in Figure 12.33. To write the equation of the normal line, let
x y z be any point on it. Then the vector
x − x0 i +
y − y0 j +
z − z0 k is along this line, hence is parallel to &
P0 . This means that, for some scalar t, either of these vectors is t times the other, say
x − x0 i +
y − y0 j +
z − z0 k = t&
P0 .

-2 z

-2 z

-1 -2

-1

-1 2 x

1

0

-2 1

-1

1

0

2 x

1 2 y

2 y

FIGURE 12.31 Part of the tangent plane to z = sin
xy at
2 1 sin
2.

FIGURE 12.32 Part of the tangent plane to z = sin
xy at
−1, −2, sin
2

 (P0 ) P0

Normal line to the surface (x, y, z)  k at P0

FIGURE 12.33

12.4 The Gradient Field and Directional Derivatives

509

The components on the left must equal the respective components on the right: x − x0 =






P0 t y − y0 =
P0 t z − z0 =
P t

x

y

z 0

These are parametric equations of the normal line. As t varies over the real line, these equations give coordinates of points
x y z on the normal line.

EXAMPLE 12.17

 Consider again the cone√
x y z = x2 + y2 − z = 0. In Example 12.15 we computed the gradient vector at
1 1 2, obtaining 1 1 − √ i − √ j + k 2 2 The normal line through this point has parametric equations √ 1 1 x − 1 = − √ t y − 1 = − √ t z − 2 = t 2 2 We can also write

SECTION 12.4

√ 1 1 x = 1 − √ t y = 1 − √ t z = 2 + t 2 2

PROBLEMS

In each of Problems 1 through 6, compute the gradient of the function and evaluate this gradient at the given point. Determine at this point the maximum and minimum rate of change of the function. 1. 
x y z = xyz
1 1 1 2. 
x y z = x2 y − sin
xz
1 −1 /4 3. 
x y z = 2xy + xe 
−2 1 6

10. 
x y z = yz + xz + xy i − 4k In each of Problems 11 through 16, find the equations of the tangent plane and normal line to the surface at the point. √ 11. x2 + y2 + z2 = 4
1 1 2 12. z = x2 + y
−1 1 2

z

4. 
x y z = cos
xyz
−1 1 /2

13. z2 = x2 − y2 
1 1 0 14. x2 − y2 + z2 = 0
1 1 0

5. 
x y z = cosh
2xy − sinh
z
0 1 1  6. 
x y z = x2 + y2 + z2 
2 2 2

15. 2x − cos
xyz = 3
1  1

In each of Problems 7 through 10, compute the directional derivative of the function in the direction of the given vector. √ 7. 
x y z = 8xy2 − xz
1/ 3
i + j + k

In each of Problems 17 through 20, find the angle between the two surfaces at the given point of intersection. (Compute this angle as the angle between the normals to the surfaces at this point).

8. 
x y z = cos
x − y + ez  i − j + 2k 9. 
x y z = x2 yz3  2¯j + k

16. 3x4 + 3y4 + 6z4 = 12
1 1 1

17. z = 3x2 + 2y2  −2x + 7y2 − z = 0
1 1 5 √ 18. x2 + y2 + z2 = 4 z2 + x2 = 2
1 2 1

510 19. z =

CHAPTER 12

Vector Differential Calculus

 √ x2 + y2  x2 + y2 = 8
2 2 8

21. Suppose & = i + k. What can be said about level surfaces of ? Prove that the streamlines of & are orthogonal to the level surfaces of .

20. x2 + y2 + 2z2 = 10 x + y + z = 5
2 2 1

12.5

Divergence and Curl The gradient operator & produces a vector field from a scalar field. We will now discuss two other vector operations. One produces a scalar field from a vector field, and the other a vector field from a vector field.

DEFINITION 12.10

Divergence

The divergence of a vector field F
x y z = f
x y zi + g
x y zj + h
x y zk is the scalar field

f g h + +  div F =

x y z

For example, if F = 2xyi +
xyz2 − sin
yzj + zex+y k, then div F = 2y + xz2 − z cos
yz + ex+y  We read div F as the divergence of F, or just “div F”.

DEFINITION 12.11

Curl

The curl of a vector field F
x y z = f
x y zi + g
x y zj + h
x y zk is the vector field      

f h

g f

h g − − − curl F = i+ j+ k

y z

z x

x y

This vector is read “curl of F”, or just “curl F”. For example, if F = yi + 2xzj + zex k, then curl F = −2xi − zex j +
2z − 1k Divergence, curl and gradient can all be thought of in terms of the vector operations of multiplication of a vector by a scalar, dot product and cross product, using the del operator &. This is defined by &=



i + j + k

x

y

z

The symbol &, which is read “del”, is treated like a vector in carrying out calculations, and the “product” of / x, / y and / z with a function 
x y z is interpreted to mean, respectively,

12.5 Divergence and Curl

511

/ x, / y and / z. In this way, the gradient of  is the product of the vector & with the scalar function :  







 i+ j+ k  = i+ j+ k

x

y

z

x

y

z = & = gradient of  The divergence of a vector is the dot product of del with the vector:  

i + j + k ·
f i + gj + hk & ·F =

x

y

z =

f g h + + = divergence of F

x y z

And the curl of a vector is the cross product of del with the vector:    i j k   & × F =  / x / y / z   f g h       

h g

f h

g f = − i+ − j+ − k

y z

z x

x y = curl F Informally, del times = gradient, del dot = divergence, and del cross = curl. This provides a way of thinking of gradient, divergence and curl in terms of familiar vector operations involving del, and will prove to be an efficient tool in carrying out computations. There are two fundamental relationships between gradient, divergence and curl. The first states that the curl of a gradient is the zero vector.

THEOREM 12.6

Curl of a Gradient

Let  be continuous with continuous first and second partial derivatives. Then & ×
& = O This conclusion can also be written curl
& = O The zero on the right is the zero vector, since the curl of a vector field is a vector.

512

CHAPTER 12 Proof

Vector Differential Calculus

By direct computation,  





 i+ j+ k & ×
& = & ×

x

y

z    i j k  

/ y

/ z  =  / x  / x / y / z   2   2   2 



2 



2 



2  = − i+ − j+ − k=O

y z z y

z x x z

x y y x

because the paired mixed partial derivatives in each set of parentheses are equal. The second relationship states that the divergence of a curl is the number zero. THEOREM 12.7

Let F be a continuous vector field whose components have continuous first and second partial derivatives. Then & ·
& × F = 0 We may also write div
curl F = 0 Proof

As with the preceding theorem, proceed by direct computation:      

h g

f h

g f − − − & ·
& × F = + +

x y z

y z x

z x y =

2 h

2 g

2 f

2 h

2 g

2 f − + − + − =0

x y x z y z y x z x z y

because equal mixed partials appear in pairs with opposite signs. Divergence and curl have physical interpretations, two of which we will now develop.

12.5.1 A Physical Interpretation of Divergence Suppose F
x y z t is the velocity of a fluid at point
x y z and time t. Time plays no role in computing divergence, but is included here because normally a velocity vector does depend on time. Imagine a small rectangular box within the fluid, as in Figure 12.34. We would like some measure of the rate per unit volume at which fluid flows out of this box across its faces, at any given time. First look at the front face II and the back face I in the diagram. The normal vector pointing out of the box from face II is i. The flux of the flow out of the box across face II is the normal component of the velocity (dot product of F with i), multiplied by the area of this face: flux outward across face II = F
x + x y z t · i y z = f
x + x y z t y z

12.5 Divergence and Curl

513

Back face I (x, y, z) z Front face II

z

(x   x, y, z)

x y

y x FIGURE 12.34

On face I, the unit outer normal is −i, so flux outward across face I = F
x y z t ·
−i y z = −f
x y z t y z The total outward flux across faces I and II is therefore

f
x + x y z t − f
x y z t y z A similar calculation can be done for the other two pairs of sides. The total flux of fluid out of the box across its faces is

f
x + x y z t −f
x y z t y z + g
x y + y z t − g
x y z t x z + h
x y z + z t − h
x y z t x y The flux per unit volume is obtained by dividing this flux by the volume x y z of the box: flux per unit volume out of the box =

f
x + x y z t − f
x y z t x g
x y + y z t − g
x y z t + y +

h
x y z + z t − h
x y z t  z

Now take the limit as x → 0, y → 0 and z → 0. The box shrinks to the point
x y z and the flux per unit volume approaches

f g h + + 

x y z which is the divergence of F
x y z t at time t. We may therefore intrepret the divergence of F as a measure of the outward flow or expansion of the fluid from this point.

12.5.2

A Physical Interpretation of Curl

Suppose an object rotates with uniform angular speed  about a line L as in Figure 12.35. The angular velocity vector  has magnitude  and is directed along L as a right-handed screw would progress if given the same sense of rotation as the object. Put L through the origin as a convenience, and let R = xi + yj + zk for any point
x y z on the rotating object. Let T
x y z be the tangential linear velocity. Then T =  R sin
 =  × R 

514

CHAPTER 12

Vector Differential Calculus (x, y, z)

T 

R θ

L

R sin θ

(0, 0, 0)

FIGURE 12.35 Angular velocity as the curl of the linear velocity.

where  is the angle between R and . Since T and  × R have the same direction and magnitude, we conclude that T =  × R. Now write  = ai + bj + ck to obtain T =  × R =
bz − cyi +
cx − azj +
ay − bxk Then

  i  & × T =  / x  bz − cy

j

/ y cx − az

k

/ z ay − bx

     

= 2ai + 2bj + 2ck = 2 Therefore 1  = & × T 2 The angular velocity of a uniformly rotating body is a constant times the curl of the linear velocity. Because of this interpretation, curl was once written rot (for rotation), particularly in British treatments of mechanics. This is also the motivation for the term irrotational for a vector field whose curl is zero. Other interpretations of divergence and curl follow from vector integral theorems we will see in the next chapter.

SECTION 12.5

PROBLEMS

In each of Problems 1 through 6, compute & · F and & × F and verify explicitly that & ·
& × F = 0.

In each of Problems 7 through 12, compute & and verify explicitly that & ×
& = O.

1. F = xi + yj + 2zk

7. 
x y z = x − y + 2z2

2. F = sinh
xyzj

8. 
x y z = 18xyz + ex

3. F = 2xyi + xey j + 2zk

9. 
x y z = −2x3 yz2

4. F = sinh
xi + cosh
xyzj −
x + y + zk

10. 
x y z = sin
xz

5. F = x2 i + y2 j + z2 k

11. 
x y z = x cos
x + y + z

6. F = sinh
x − zi + 2yj +
z − y2 k

12. 
x y z = ex+y+z

12.5 Divergence and Curl 13. Let 
x y z be a scalar field and F a vector field. Derive expressions for & ·
F and & ×
F in terms of operations applied to 
x y z and F. 14. Let F = f i + gj + hk be a vector field. Define      

i+ g j+ h k F·& = f

x

y

z

515

Let G be a vector field. Show that &
F·G =
F·&G+
G·&F+F×
& ×G+G×
& ×F 15. Let F and G be vector fields. Prove that & ·
F × G = G ·
& × G − F ·
& × G 16. Let 
x y z and 
x y z be scalar fields. Prove that & ·
& × & = 0

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CHAPTER

13

LINE INTEGRALS GREEN’S THEOREM INDEPENDENCE OF PATH AND POTENTIAL THEORY IN THE PLANE SURFACES IN 3-SPACE AND SURFACE INTEGRALS APPLICATIONS OF SURFACE INTEGRALS PREPARA

Vector Integral Calculus

This chapter is devoted to integrals of vector fields over curves and surfaces, and relationships between such integrals. These have important uses in solving partial differential equations and in constructing models used in the sciences and engineering.

13.1

Line Integrals We begin with the integral of a vector field over a curve. This requires some background on curves. Suppose a curve C in 3-space is given by parametric equations x = x
t y = y
t z = z
t

for a ≤ t ≤ b

We call x
t, y
t and z
t the coordinate functions of C. We will think of C not only as a geometric locus of points
x
t y
t z
t, but also as having an orientation or direction, given by the direction this point moves along C as t increases from a to b. Denote this orientation by putting arrows on the graph of the curve (Figure 13.1). Trajectories of a system of differential equations are also oriented curves, with the particle moving along the geometric locus in a direction dictated by the flow of the system. We call
x
a y
a z
a the initial point of C, and
x
b y
b z
b the terminal point. A curve is closed if the initial and terminal points are the same.

517

518

CHAPTER 13

Vector Integral Calculus z

z

(x(b), y(b), z(b)) (2, 0, 4)

R

y

y x (x(a), y(a), z(a)) 0

x FIGURE 13.1

Orientation along a

2

t

FIGURE 13.2

curve. EXAMPLE 13.1

Let C be given by x = 2 cos
t y = 2 sin
t z = 4

0 ≤ t ≤ 2.

A graph of C is shown in Figure 13.2. This graph is the circle of radius 2 about the origin in the plane z = 4. The arrow on the curve indicates its orientation (the direction of motion of
2 cos
t 2 sin
t 4 around the graph as t varies from 0 to 2). The initial point is
2 0 4, obtained at t = 0, and the terminal point is also
2 0 4, obtained at t = 2. This curve is closed. Contrast C with the curve K given by x = 2 cos
t y = 2 sin
t z = 4 0 ≤ t ≤ 3. K has the same graph and counterclockwise orientation as C, but is a different curve. K is not closed. The initial point of K is
2 0 4, but its terminal point is
−2 0 4. A particle moving about K goes around the circle once, and then makes another half-circle. It would presumably take more energy to move a particle about K than about C. K is also longer than C. A curve is: – continuous if its coordinate functions are continuous on the parameter interval, – differentiable if its coordinate functions are differentiable, and – smooth if its coordinate functions have continuous derivatives, which are not all zero for the same value of t. Because x

ti + y

tj + z

tk is tangent to the curve if this is not the zero vector, a smooth curve is one that has a continuous tangent vector. The graph of a curve consists of all points
x
t y
t z
t as t varies over the parameter interval. As Example 13.1 shows, there is more to a curve than just its graph. We are dealing with curves both as geometric objects and as oriented objects, having a sense of direction along the graph. We will now define the line integral of a vector field over a smooth curve. DEFINITION 13.1

Line Integral

Suppose a smooth curve C has coordinate functions x = x
t, y = y
t, z = z
t for a ≤ t ≤ b. Let f
x y z, g
x y z, and h
x y z be continuous at least on the graph of C. Then, the line integral
f
x y z dx + g
x y z dy + h
x y z dz C

13.1 Line Integrals

519

is defined by
f
x y z dx + g
x y z dy + h
x y z dz C

=


b a

dy dz dx f
x
t y
t z
t + g
x
t y
t z
t + h
x
t y
t z
t dt dt dt dt
131

We can write this line integral more compactly as
C

fdx + gdy + hdz

 To evaluate C fdx + gdy + hdz, substitute the coordinate functions x = x
t, y = y
t and z = z
t into f
x y z, g
x y z and h
x y z, obtaining functions of t. Further, substitute dx =

dx dt, dt

dy =

dy dt dt

and dz =

dz dt dt

This results in the Riemann integral on the right side of equation (13.1) of a function of t over the range of values of this parameter.

EXAMPLE 13.2

Evaluate the line integral

 C

xdx − yzdy + ez dz if C is given by x = t3  y = −t z = t2  1 ≤ t ≤ 2

Here f
x y z = x

g
x y z = −yz

and h
x y z = ez

and, on C, dx = 3t2 dt

dy = −dt

and dz = 2tdt.

Then
C

= =

xdx − yzdy + ez dz


2 1


2

 2 t3
3t2  −
−t
t2 
−1 + et
2t dt

1

 111 2 3t5 − t3 + 2tet dt = + e4 − e 4

520

CHAPTER 13

Vector Integral Calculus

EXAMPLE 13.3

 Evaluate C xyzdx −cos
yzdy +xzdz over the straight line segment from
1 1 1 to
−2 1 3. Here we are left to find the coordinate functions of the curve. Parametric equations of the line through these points are x
t = 1 − 3t

y
t = 1

z
t = 1 + 2t

We must let t vary from 0 to 1 for the initial point to be
1 1 1 and the terminal point to be
−2 1 3. Now
xyzdx − cos
yzdy + xzdz C

= =




1 0


1 − 3t
1
1 + 2t
−3 − cos
1 + 2t
0 +
1 − 3t
1 + 2t
2 dt


1 0

 3 −1 + t + 6t2 dt = . 2

If C is a smooth curve in the x y plane (zero z-component), and f
x y and g
x y are continuous on C, then we can write a line integral
f
x ydx + g
x ydy C

which we refer to as a line integral in the plane. We evaluate this according to the Definition 13.1, except now there is no z-component.

EXAMPLE 13.4

 Evaluate C xydx − y sin
xdy if C is given by x
t = t2 and y
t = t for −1 ≤ t ≤ 4. Proceed:
xydx − y sin
xdy C

=




4 −1

t2 t
2t − t sin
t2 
1dt

= 410 +

1 1 cos
16 − cos
1 2 2

In these examples we have included all the terms to follow equation 13.1 very literally. However, with some experience there are obvious shortcuts one can take. In Example 13.3, for could have example, y is constant on the curve, so dy = 0 and the term g
x
t y
t z
t dy dt been simply omitted. Thus far we can integrate only over a smooth curve. This requirement can be relaxed as follows. A curve C is piecewise smooth if x

t, y

t and z

t are continuous, and not all zero for the same value of t, at all but possibly finitely many values of t. Since x

ti + y

tj + z

tk is the tangent vector to C if this is not the zero vector, this condition means that a piecewise smooth curve has a continuous tangent at all but finitely many points. Such a curve typically has the appearance of Figure 13.3, with smooth pieces C1      Cn connected at points where the curve may have no tangent. We will refer to a piecewise smooth curve as a path. In Figure 13.3 the terminal point of Cj is the initial point of Cj+1 for j = 1     n − 1. The segments C1      Cn are in order as one moves from the initial to the terminal point of C. This

13.1 Line Integrals z

C4

y

C3 C1

C2

521

(0, 1)

C2

(2, 1)

C1

y

(1, 0)

x FIGURE 13.3 Typical piecewise smooth curve (path).

x

FIGURE 13.4

is indicated by the arrows showing orientation along the smooth pieces of C. If f , g and h are continuous over each Cj , then we define


fdx + gdy + hdz = fdx + gdy + hdz + · · · + fdx + gdy + hdz C

C1

Cn

This allows us to take line integrals over paths, rather than restricting the integral to smooth curves.

EXAMPLE 13.5

Let C be the curve consisting of the quarter circle x2 + y2 = 1 in the x y− plane, from
1 to
0 1, followed by the horizontal line segment from
0 1 to
2 1. Compute  0 2 x ydx + y2 dy. C C is piecewise smooth and consists of two smooth pieces (Figure 13.4). Parametrize these as follows. The quarter circle part is C1  x = cos
t y = sin
t

0 ≤ t ≤ /2

The straight segment is C2  x = p y = p

0 ≤ p ≤ 2

We have used different names for the parameters on the curves to help distinguish them. Now evaluate the line integral along each of the two curves making up the path. For the line integral over C1 , compute

/2 x2 ydx + y2 dy =

cos2
t sin
t
− sin
t + sin2
t cos
tdt 0

C1

=


0

/2 

 1 1 − cos2
t sin2
t + sin2
t cos
t dt = −  +  16 3

Next evaluate the line integral over C2 . On C2 , x = p and y = 1, so dy = 0 and

2 8 x2 ydx + y2 dy = p2 dp =  3 C2 0 Then
C

x2 ydx + y2 dy = −

1 8  1  + + = − + 3 16 3 3 16

522

CHAPTER 13

Vector Integral Calculus

It is sometimes useful to think of a line integral in terms of vector  operations, particularly in the next section when we deal with potential functions. Consider C fdx + gdy + hdz. Form a vector field F
x y z = f
x y zi + g
x y zj + h
x y zk If C has coordinate functions x = x
t, y = y
t, z = z
t, we can form the position vector, for C: R
t = x
ti + y
tj + z
tk At any time t, the vector R
t can be represented by an arrow from the origin to the point
x
t y
t z
t on C. As t varies, this vector pivots about the origin and adjusts its length to sweep out the curve. If C is smooth, then the tangent vector R

t is continuous. Now dR = dxi + dyj + dzk, so F · dR = f
x y zdx + g
x y zdy + h
x y zdy and
C

f
x y zdx + g
x y zdy + h
x y zdz =


C

F · dR

This is just another way of writing a line integral in terms of vector operations. Line integrals arise in many contexts. For example, consider a force F causing a particle to move along a smooth curve C having position function R
t, where t varies from a to b. At any point
x
t y
t z
t of C, the particle may be thought of as moving in the direction of the tangent to its trajectory, and this tangent is R

t. Now F
x
t y
t z
t · R

t is the dot product of a force with a direction, and so has the dimensions of work. By integrating this function from a to b, we “sum” the work being done by F over the entire path of motion. This  suggests that C F · dR, or C fdx + gdy + hdz, can be interpreted as the work done by F in moving the particle over the path.

EXAMPLE 13.6

Calculate the work done by F = i − yj + xyzk in moving a particle from
0 0 0 to
1 −1 1 along the curve x = t, y = −t2 , z = t for 0 ≤ t ≤ 1. The work is

work = F · dR = dx − ydy + xyzdz C

= =

C


1 0

 1 + t2
−2t − t4 dt


1 0

 3 1 − 2t3 − t4 dt =  10

The correct units (such as foot-pounds) would have to be provided from context. Line integrals have some of the usual properties we associate with integrals.

13.1 Line Integrals

523

THEOREM 13.1

Let C be a path and let having position vector R. Let F and G be vector fields that are continuous at points of C. Then 1.


C


F + G · dR =

2. For any number ,


C


C

F · dR +

F · dR = 


C


C

G · dR

F · dR

This theorem illustrates the efficiency of the vector notation for line integrals. We could also write the conclusion (1) as

f + f ∗ dx +
g + g ∗ dy +
h + h∗ dz C

=




C

fdx + gdy + hdz +


C

f ∗ dx + g ∗ dy + h∗ dz

 a integrals, reversing the limits of integration changes the sign of the integral:  b For Riemann f
xdx = − f
xdx. The analogous property for line integrals involves reversing orientaa b tion on C. Given C with an orientation from an initial point P to a terminal point Q, let −C denote the curve obtained from C by reversing the orientation to go from Q to P (Figure 13.5). Here is a more careful definition of this orientation reversal. z

z

C


C

Q y

P

Q y

P

x

x

FIGURE 13.5 Reversing orientation of a curve.

DEFINITION 13.2

Let C be a smooth curve with coordinate functions x = x
t, y = y
t, z = z
t, for a ≤ t ≤ b. Then −C denotes the curve having coordinate functions , x
t = x
a + b − t

, y
t = y
a + b − t , z
t = z
a + b − t

for a ≤ t ≤ b.

The initial point of −C is
, x
a, y
a, z
a =
x
b y
b z
b

524

CHAPTER 13

Vector Integral Calculus

the terminal point of C. And the terminal point of −C is
, x
b, y
b, z
b =
x
a y
a z
a the initial point of C. By the chain rule, −C is piecewise smooth if C is piecewise smooth. We will now show that the line integral of a vector field over −C is the negative of the line integral of the vector field over C. THEOREM 13.2

Let C be a smooth curve with coordinate functions x = x
t, y = y
t, z = z
t. Let f , g and h be continuous on C. Then

fdx + gdy + hdz = − fdx + gdy + hdz −C

First,


Proof

C

= Similarly,

fdx + gdy + hdz


b dx dy dz f
x
t y
t z
t + g
x
t y
t z
t + h
x
t y
t z
t dt dt dt dt a
−C

=

C

fdx + gdy + hdz


b a

d, y d, z d, x x
t, y
t, z
t + h
, x
t, y
t, z
t f
, x
t, y
t, z
t + g
, dt dt dt dt

Change variables in the last integral by putting s = a + b − t. When t = a, s = b , and when t = b, s = a. Further, d, x d d dx ds dx = x
a + b − t = x
s = =− dt dt dt ds dt ds and, similarly, d, y dy d, z dz =−  =−  dt ds dt ds Finally, dt = −ds Then


−C

fdx + gdy + hdz

dy dz dx − g
x
s y
s z
s − h
x
s y
s z
s
−1ds ds ds ds b


b dy dz dx f
x
s y
s z
x + g
x
s y
s z
s + h
x
s y
s z
s =− ds ds ds ds a
= − fdx + gdy + hdz

=




a

−f
x
s y
s z
x

C

13.1 Line Integrals

525

 In view of this  theorem, the easiest way to evaluate −C fdx + gdy + hdz is usually to take the negative of C fdx + gdy + hdz. We need not actually write the coordinate functions of −C, as was done in the proof.

EXAMPLE 13.7

A force F
x y z = x2 i − zyj + x cos
zk moves a particle along the path C given by x = t2 , y = t, z = t for 0 ≤ t ≤ 3. The initial point is P 
0 0 0 and the terminal point of C is Q 
9 3 3. Suppose we want the work done in moving the particle along this path from Q to P.  Since we want to go from the terminal to the initial point of C, the work done is −C F · dR. However, we do not need to formally define −C in terms of new coordinate functions. We can simply calculate C F · dR, and take the negative of this. Calculate

F · dR = fdx + gdy + hdz C

C

=


3

 t4
2t − t
t
1 + t2 cos
t
 dt

0

6   The work done in moving the particle along the path from Q to P is = 243 − 9 −

6 + 9 − 243 

13.1.1

Line Integral with Respect to Arc Length

Line integrals with respect to arc length occur in some uses of line integrals. Here is the definition of this kind of line integral.

DEFINITION 13.3

Line Integral With Respect to Arc Length

Let C be a smooth curve with coordinate functions x = x
t, y = y
t, z = z
t for a ≤ t ≤ b. Let  be a real-valued function that is continuous on the graph of C. Then the integral of  over C with respect to arc length is
b
 
x y zds = 
x
t y
t z
t x

t2 + y

t2 + z

t2 dt C

a

The rationale behind this definition is that the length function along C is
t s
t = x

2 + y

2 + z

2 d a

Then ds =



x

t2 + y

t2 + z

t2 dt

suggesting the integral in the definition.

526

CHAPTER 13

EXAMPLE 13.8

Evaluate

Vector Integral Calculus

 C

xyds over the curve given by x = 4 cos
t y = 4 sin
t z = −3

for 0 ≤ t ≤ /2

Compute
C

xyds = =




/2 0




/2

0

" 4 cos
t 4 sin
t 16 sin2
t + 16 cos2
tdt 64 cos
t sin
tdt = 32

Line integrals with respect to arc length occur in calculations of mass, density, and various other quantities for one-dimensional objects. Suppose, for example, we want the mass of a thin wire bent into the shape of a piecewise smooth curve C having coordinate functions x = x
t y = y
t z = z
t

for a ≤ t ≤ b

The wire is one-dimensional in the sense that (ideally) it has length but not area or volume. We will derive an expression for the mass of the wire as follows. Let 
x y z be the density of the wire at any point. Partition a b into subintervals by inserting points a = t0 < t1 < · · · < tn−1 < tn = b Choose these points t units apart, so tj − tj−1 = t. These partition points of a b determine points Pj 
x
tj  y
tj  z
tj  along C, as shown in Figure 13.6. Assuming that the density function is continuous, we can choose t sufficiently small that on the piece of wire between Pj−1 and Pj , the values of the density function are approximated to whatever accuracy we wish by 
Pj . The length of the segment of wire between Pj−1 and Pj is s = s
Pj  − s
Pj−1 , which is approximated by ds =

"

x

tj 2 +
y

tj 2 + z

tj 2 t

The density of this piece of wire between Pj−1 and Pj is therefore approximately the “nearly” constant value of the density on this piece, times the length of this piece of wire, this product being " 
x
tj  y
tj  z
tj  x

tj 2 +
y

tj 2 + z

tj 2 t

z Pj
1

Pn : (x(b), y(b), z(b))

Pj P0: (x (a), y(a), z(a)) x FIGURE 13.6

t0  a

y t1 t2 t3

b  tn

t

13.1 Line Integrals

527

The mass of the entire length of wire is approximately the sum of the masses of these pieces: " n  mass ≈ 
x
tj  y
tj  z
tj  x

tj 2 +
y

tj 2 + z

tj 2 t j=1

in which ≈ means “approximately equal”. Recognize this as the Riemann sum for a definite integral to obtain, in the limit as t → 0,
 mass = 
x
t y
t z
t x

t2 + y

t2 + z

t2 dt C

=





x y zds C

A similar argument leads to coordinates
x y z of the center of mass of the wire: 1
1
1
x= x
x y zds y = y
x y zds z = z
x y zds m C m C m C in which m is the mass of the wire.

EXAMPLE 13.9

A wire is bent into the shape of the quarter circle C given by x = 2 cos
t y = 2 sin
t z = 3

for 0 ≤ t ≤ /2

The density function is 
x y z = xy2 grams/centimeter. We want the mass and center of mass of the wire. The mass is
m = xy2 ds = =





C

/2

0 /2 0

" 2 cos
t 2 sin
t2 4 sin2
t + 4 cos2
tdt 16 cos
t sin2
tdt =

16 grams. 3

Now compute the coordinates of the center of mass. First, 1
x= x
x y zds m C " 3
/2 =

2 cos
t2 2 sin
t2 4 sin2
t + 4 cos2
tdt 16 0
/2 3 =6  cos2
t sin2
tdt = 8 0 Next,

1
y
x y zds m C " 3
/2 =

2 cos
t 2 sin
t3 4 sin2
t + 4 cos2
tdt 16 0
/2 3 =6 cos
t sin3
tdt =  2 0

y=

CHAPTER 13

528

Vector Integral Calculus

Finally, 1
z
x y zds m C " 3
/2 2 = 3 2 cos
t 2 sin
t 4 sin2
t + 4 cos2
tdt 16 0
/2 =9 sin2
t cos
tdt = 3

z=

0

The last result could have been anticipated, since the z-component on the curve is constant. The center of mass is
3/8 3/2 3 

SECTION 13.1

PROBLEMS

In each of Problems 1 through 15, evaluate the line integral.  1. C xdx − dy + zdz, with C given by x
t = t y
t = t z
t = t3 for 0 ≤ t ≤ 1  2. C −4xdx + y2 dy − yzdz, with C given by x
t = −t2  y
t = 0 z
t = −3t for 0 ≤ t ≤ 1  3. C
x + yds, where C is given by x = y = t z = t2 for 0≤t≤2  4. C x2 zds, where C is the line segment from
0 1 1 to
1 2 −1  5. C F · dR, where F = cos
xi − yj + xzk and R = ti − t2 j + k for 0 ≤ t ≤ 3  6. C 4xyds, with C given by x = y = t z = 2t for 1≤t≤2  7. C F · dR, with F = xi + yj − zk and C the circle x2 + y2 = 4 z = 0, going around once counterclockwise.  8. C yzds, with C the parabola z = y2  x = 1 for 0 ≤ y≤2  √ 9. C −xyzdz, with C the curve x = 1 y = z for 4 ≤ z≤9

13.2

10. 11. 12. 13. 14. 15.



xzdy, with C the curve x = y = t z = −4t2 for 1≤t≤3  8z2 ds, with C the curve x = y = 2t2  z = 1 for C 1≤t≤2  F · dR, with F = i − xj + k and R = cos
ti − C sin
tj + tk for 0 ≤ t ≤   8x2 dy, with C given by x = et  y = −t2  z = t for C 1≤t≤2  xdy − ydx, C the curve x = y = 2t z = e−t for C 0≤t≤3  sin
xds, with C given by x = t y = 2t z = 3t for C 1≤t≤3 C

16. Find the mass and center of mass of a thin, straight wire from the origin to
3 3 3 if 
x y z = x + y + z grams per centimeter. 17. Find the work done by F = x2 i − 2yzj + zk in moving an object along the line segment from
1 1 1 to
4 4 4. b 18. Show that any Riemann integral a f
xdx is equal  to a line integral C F · dR for appropriate choices of F and C. In this sense the line integral generalizes the Riemann integral.

Green’s Theorem Green’s theorem was developed independently by the self-taught British amateur natural philosopher George Green and the Ukrainian mathematician Michel Ostrogradsky. They were studying potential theory (electric potentials, potential functions), and they obtained an important relationship between double integrals and line integrals in the plane.

13.2 Green’s Theorem

529

y Exterior of C

y

Interior of C

C x

x FIGURE 13.7 Graph of a curve that is not simple.

FIGURE 13.8

The Jordan

curve theorem.

Let C be a piecewise smooth curve in the plane, having coordinate functions x = x
t, y = y
t for a ≤ t ≤ b. We will be interested in this section in C being a closed curve, so the initial and terminal points coincide. C is positively oriented if
x
t y
t moves around C counterclockwise as t varies from a to b. If
x
t y
t moves clockwise, then we say that C is negatively oriented. For example, let x
t = cos
t and y
t = sin
t for 0 ≤ t ≤ 2. Then
x
t y
t moves counterclockwise once around the unit circle as t varies from 0 to 2, so C is positively oriented. If, however, K has coordinate functions x
t = − cos
t and y
t = sin
t for 0 ≤ t ≤ 2, then K is negatively oriented, because now
x
t y
t moves in a clockwise sense. However, C and K have the same graph. A closed curve in the plane is positively oriented if, as you walk around it, the region it encloses is over your left shoulder. A curve is simple if the same point cannot be on the graph for different values of the parameter. This means that x
t1  = x
t2  and y
t1  = y
t2  can occur only if t1 = t2 . If we envision the graph of a curve as a train track, this means that the train does not return to the same location at a later time. Figure 13.7 shows the graph of a curve that is not simple. This would prevent a closed curve from being simple, but we make an exception of the initial and terminal points. If these are the only points obtained for different values of the parameter, then a closed curve is also called simple. For example, the equations x = cos
t and y = sin
t for 0 ≤ t ≤ 2, describe a simple closed curve. However, consider M given by x = cos
t, y = sin
t for 0 ≤ t ≤ 4. This is a closed curve, beginning and ending at
1 0, but
x
t y
t traverses the unit circle twice counterclockwise as t varies from 0 to 4. M is a closed curve but it is not simple. It is a subtle theorem of topology, the Jordan curve theorem, that a simple closed curve C in the plane separates the plane into two regions having C as common boundary. One region contains points arbitrarily far from the origin, and is called the exterior of C. The other region is called the interior of C. These regions are displayed for a typical closed curve in Figure 13.8. The interior of C has finite area, while the exterior does not. / Finally, when a line integral is taken around a closed curve, we often use the symbol C in place of C . This notation is optional, and is simply a reminder that C is closed. It does not alter in any way the meaning of the integral. We are now ready to state the first fundamental theorem of vector integral calculus. Recall that a path is a piecewise smooth curve (having a continuous tangent at all but finitely many points).

THEOREM 13.3

Green

Let C be a simple closed positively oriented path in the plane. Let D consist of all points on C and in its interior. Let f , g, g/ x and f/ y be continuous on D. Then

 g f  − f
x ydx + g
x ydy = dA

x y C D

530

CHAPTER 13

Vector Integral Calculus

The significance of Green’s theorem is that it relates an object that deals with a curve, which is one-dimensional, to an object related to a planar region, which is two-dimensional. This will have important implications when we discuss independence of path of line integrals in the next section, and later when we develop partial differential equations and complex analysis. Green’s theorem will also lead shortly to Stokes’s theorem and Gauss’s theorem, which are its generalizations to 3-space. We will prove the theorem under restricted conditions at the end of this section. For now, here are two computational examples.

EXAMPLE 13.10

Sometimes we use Green’s theorem as a computational aid to convert one kind of integral into another, possibly simpler, one. As an illustration, suppose we want to compute the work done by the force F
x y =
y − x2 ex i +
cos
2y2  − xj in moving a particle about the / rectangular path C of Figure 13.9. If you try to evaluate C F · dR as a sum of line integrals over the straight line sides of this rectangle, you will find that the integrations cannot be done in elementary form. However, apply Green’s theorem, with D the region bounded by the rectangle. We obtain 



2 2 x work = F · dR =
cos
2y  − x −
y − x e  dA

x

y C D

= −2dA =
−2(area of D) = −4 D

EXAMPLE 13.11

Another typical use of Green’s theorem is in deriving very general results. To illustrate, suppose we want to evaluate 2x cos
2ydx − 2x2 sin
2ydy C

for every positively oriented simple closed path C in the plane. This may appear to be a daunting task, since there are infinitely many different such paths. However, observe the form of f
x y and g
x y in the line integral. In particular, 

 −2x2 sin
2y −
2x cos
2y = −4x sin
2y + 4x sin
2y = 0

x

y for all x and y. Therefore, Green’s theorem gives us

2x cos
2ydx − 2x2 sin
2ydy = 0dA = 0 C

D

In the next section we will see how the vanishing of this line integral for any closed curve allows an important conclusion about line integrals K 2x cos
2ydx − 2x2 sin
2ydy when K is not closed. We will conclude this section with a proof of Green’s theorem under special conditions on the region bounded by C. Assume that D can be described in two ways.

13.2 Green’s Theorem

531

y (0, 3) D

C2

y

(1, 3)

y  k(x), x from b to a

C (1, 1)

C1

(0, 1) x

a

FIGURE 13.9

y  h(x), x from a to b x b

FIGURE 13.10

First, D consists of all points
x y with a≤x≤b and, for each x, h
x ≤ y ≤ k
x Graphs of the curves y = h
x and y = k
x form, respectively, the lower and upper parts of the boundary of D (see Figure 13.10). Second, D consists of all points
x y with c≤y≤d and, for each y, F
y ≤ x ≤ G
y In this description, the graphs of x = F
y and x = G
y form, respectively, the left and right parts of the boundary of D (see Figure 13.11). Using these descriptions of D and the boundary of D, we can demonstrate Green’s theorem by evaluating the integrals involved. First, let C1 be the lower part of C (graph of y = h
x) and C2 the upper part (graph of y = k
x). Then

f
x ydx = f
x ydx + f
x ydx C

C1

= =




b

a




b

a

C2

f
x h
xdx +




a

f
x k
xdx b

− f
x k
x − f
x h
xdx

The upper and lower limits of integration in the second line maintain a counterclockwise orientation on C. y y d x  F( y)

x  G(y) x

y c FIGURE 13.11

532

CHAPTER 13

Vector Integral Calculus

Now compute
b
k
x f

f dA = dydx a h
x y D y
b k
x =

f
x yh
x dx a

= Therefore

C




a

b

f
x k
x − f
x h
xdx

f
x ydx = −



f dA D y

Using the other description of D, a similar computation shows that

g g
x ydy = dA C D x Upon adding the last two equations we obtain the conclusion of Green’s theorem.

SECTION 13.2

PROBLEMS

1. A particle moves once counterclockwise about the triangle with vertices
0 0
4 0 and
1 6, under the influence of the force F = xyi + xj. Calculate the work done by this force. 2. A particle moves once counterclockwise about the circle of radius 6 about the origin, under the influence of the force F =
ex − y + x cosh
xi +
y3/2 + xj. Calculate the work done. 3. A particle moves once counterclockwise about the rectangle with vertices
1 1,
1 7,
3 1 and
3 7, under the influence of the force F =
− cosh
4x4  + xyi +
e−y + xj. Calculate the work done. In each of/ Problems 4 through 11, use Green’s theorem to evaluate C F · dR. All curves are oriented counterclockwise. 4. F = 2yi − xj, C is the circle of radius 4 about
1 3 5. F = x2 i − 2xyj, C is the triangle with vertices
1 1
4 1
2 6 6. F =
x + yi +
x − yj, C is the ellipse x2 + 4y2 = 1 7. F = 8xy2 j, C is the circle of radius 4 about the origin

8. F =
x2 − yi +
cos
2y − e3y + 4xj, with C any square with sides of length 5 9. F = ex cos
yi − ex sin
yj, C is any simple closed piecewise smooth curve in the plane 10. F = x2 yi − xy2 j, C the boundary of the region x2 + y2 ≤ 4, x ≥ 0, y ≥ 0 11. F = xyi +
xy2 − ecos
y j, C the triangle with vertices
0 0
3 0
0 5 12. Let C be a positively oriented simple closed path with interior D. / (a) Show that the area of D equals C −ydx. / (b) Show that the area of D equals C xdy. 1−ydx + xdy. (c) Show that the area of D equals 2 13. Let u
x y be continuous with continuous first and second partial derivatives on a simple closed path C and throughout the interior D of C. Show that

2 u 2 u

u

u − dx + dy = + 2 dA 2

y

x

y C D x

13.2.1 An Extension of Green’s Theorem There is an extension of Green’s theorem to include the case that there are finitely many points enclosed by C at which f , g, f/ y and/or g/ x are not continuous, or perhaps are not even defined. The idea is to excise the “bad points,” as we will now describe.

13.2 Green’s Theorem

533

Suppose C is a simple closed positively oriented path in the plane enclosing a region D. Suppose f , g, f/ y and g/ x are continuous on C, and throughout D except at points P1      Pn . Green’s theorem does not apply to this region. But with a little imagination we can still draw an interesting conclusion. Enclose each Pj with a circle Kj of sufficiently small radius that none of these circles intersects either C or each other (Figure 13.12). Next, cut a channel in D from C to K1 , then from K1 to K2 , and so on until finally a channel is cut from Kn−1 to Kn . A typical case is shown in Figure 13.13. Form the closed path C ∗ consisting of C (with a small segment cut out where the channel to K1 was made), each of the Kj
s (with small cuts removed where the channels entered and exited), and the segments forming the connections between C and the successive Kj
s. Figure 13.14 shows C ∗ , which encloses the region D∗ . By the way C ∗ was formed, the points P1      Pn are external to C ∗ (Figure 13.15). Further, f , g, f/ y and g/ x are continuous on C ∗ and throughout D∗ . We can therefore apply Green’s theorem to C ∗ and D∗ to conclude that

 g f  − dA f
x ydx + g
x ydy =

x y C∗ D∗

-

(13.2)

Now imagine that the channels that were cut become narrower, merging to form segments between C and successive Kj
s. Then C ∗ approaches the curve Cˆ of Figure 13.16, and D∗ ˆ shown in Figure 13.17. D ˆ consists of D with the disks bounded by approaches the region D K1      Kn cut out. In this limit process, equation 13.2 approaches C

f
x ydx + g
x ydy +



 g f  − f
x ydx + g
x ydy = dA ˆ

x y Kj D

n  j=1

C

P1 K1

C P1

P2

Kn

P2

K1

K2

K2 Kn

Pn

FIGURE 13.12

Pn

FIGURE 13.13

y P1

C*

P2

P1

P2

P3 Pn

Pn x

FIGURE 13.14

FIGURE 13.15

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C

D

FIGURE 13.16

FIGURE 13.17

On the left side of this equation, line integrals over the internal segments connecting C and the Kj
s cancel because the integration is carried out twice over each segment, once in each direction. Further, the orientation on C is counterclockwise, but the orientation on each Kj is clockwise because of the way the boundaries were traversed (Figure 13.14). If we reverse the orientations on these circles the line integrals over them change sign and we can write C

f
x ydx + g
x ydy =



 g f  − f
x ydx + g
x ydy + dA ˆ

x y D Kj

n  j=1

in which all the integrals are in the positive, counterclockwise sense about the curves C and K1      Kn . This is the generalization of Green’s theorem that we sought.

EXAMPLE 13.12

Suppose we are interested in -

−y C

x2 + y 2

dx +

x x2 + y 2

dy

with C any simple closed positively oriented path in the plane, but not passing through the origin. With f
x y =

−y x2 + y 2

and

g
x y =

x x2 + y 2



we have y 2 − x2

g

f = 2  =

x x + y2

y f , g, f/ y and g/ x are continuous at every point of the plane except the origin. This leads us to consider two cases. Case 1—C does not enclose the origin. Now Green’s theorem applies and C



 g f  −y x dA = 0 dx + 2 dy = − x2 + y 2 x + y2

x y D

13.2 Green’s Theorem y

535

C K x

FIGURE 13.18

Case 2—C encloses the origin. Draw a circle K centered at the origin, with radius sufficiently small that K does not intersect C (Figure 13.18). By the extension of Green’s theorem, f
x ydx + g
x ydy C



 g f  − dA = f
x ydx + g
x ydy + ˆ

x y K D = f
x ydx + g
x ydy -

K

ˆ is the region between K and C. Both of these line integrals are in the counterclockwise where D sense about the respective curves. The last line integral can be evaluated explicitly because we know K. Parametrize K by x = r cos
 y = r sin
 for 0 ≤  ≤ 2 Then

-

= =




K

f
x ydx + g
x ydy

2 0


0

2



 −r sin
 r cos


−r sin
 +

r cos
 d r2 r2

d = 2

We conclude that  0 f
x ydx + g
x ydy = C 2

SECTION 13.2

x y i+ 2 j x2 + y 2 x + y2  3/2 1
xi + yj 2. F = x2 + y 2

if C encloses the origin



PROBLEMS

/ In each of Problems 1 through 5, evaluate C F · dR over any simple closed path in the x y plane that does not pass through the origin. 1. F =

if C does not enclose the origin



   −y x 2 i + + x − 2y j x2 + y 2 x2 + y 2     −y x 4. F = + 3x i + 2 −y j x2 + y 2 x + y2     x y 2 5. F =  + 2x i +  − 3y j x2 + y 2 x2 + y 2 3. F =

536

13.3

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Independence of Path and Potential Theory in the Plane In physics, a conservative force field is one that is derivable from a potential. We will use the same terminology.

DEFINITION 13.4

Conservative Vector Field

Let D be a set of points in the plane. A vector field F
x y is conservative on D if for some real-valued 
x y, F = & for all
x y in D. In this event,  is a potential function for F on D.

If  is a potential function for F, then so is  + c for any constant c, because &
 + c = &. For this reason we often speak of a potential function for F, rather than the potential function. Recall that, if F
x y = f
x yi + g
x yj and R
t = x
ti + y
tj is a position function for C, then C F · dR is another way of writing C f
x ydx + g
x ydy. We will make frequent use of the notation C F · dR throughout this section because we want to examine the effect on this integral when F has a potential function, and for this we will use vector notation. First, the line integral of a conservative vector field can be evaluated directly in terms of a potential function. For suppose C is smooth, with coordinate functions x = x
t, y = y
t for a ≤ t ≤ b. If F = &, then F
x y = and


C



 i+ j

x

y




 dx + dy

y C x
b   dx  dy  = + dt

x dt y dt a
b d 
x
t y
tdt = a dt

F · dR =

= 
x
b y
b − 
x
a y
a Denoting P1 =
x
b y
b and P0 =
x
a y
a, this result states that
F · dR = 
P1  − 
P0  C

= 
terminal point of C − 
initial point of C

(13.3)

The line integral of a conservative vector field over a path is the difference in values of a potential function at end points of the path. This is familiar from physics. If a particle moves along a path under the influence of a conservative force field, then the work done is equal to the difference in the potential energy at the ends of the path. One ramification of equation (13.3) is that the actual path itself does not influence the outcome, only the end points of the path. If we chose a different path K between the same end points, we would obtain the same result for K F ·dR. This suggests the concept of independence of path of a line integral.

13.3 Independence of Path and Potential Theory in the Plane

DEFINITION 13.5



537

Independence of Path

F · dR is independent of path on a set D of points in the plane if for any points P0 and C P1 in D, the line integral has the same value over any paths in D having initial point P0 and terminal point P1 .

The discussion preceding the definition may now be summarized. THEOREM 13.4

Let  and its first partial derivatives be continuous for all
x y in a set D of points in the plane. Let F = &. / Then C F · dR is independent of path in D. Further, if C is a simple closed path in D, then C F · dR = 0. The independence of path follows from equation (13.3), which states that, when F = &,  the value of C F · dR depends only on the values of 
x y at the end points of the path, and not where the path goes in between. For the last conclusion of the theorem, if C is a closed path in D, then the initial and terminal points coincide, hence the difference between the values of  at the terminal and initial points is zero.

EXAMPLE 13.13

Let F
x y = 2x cos
2yi − 2x2 sin
2yj. It is routine to check that 
x y = x2 cos
2y is a potential function for F. Since  is continuous with continuous partial derivatives over the entire plane, we can let D consist of all points in the plane in the definition of independence of path. For example, if C is any path in the plane from
0 0 to
1 /8, then √
2  F · dR = 
1 /8 − 
0 0 = 2 C / Further, if K is any simple closed path in D, then K F · dR = 0. It is clearly to our advantage to know whether a vector field is conservative, and, if it is, to be able to produce a potential function. Let F
x y = f
x yi + g
x yj F is conservative exactly when, for some , F = & =



 i + j

x

y

and this requires that



 = f
x y and = g
x y

x

y To attempt to find such a , begin with either of these equations and integrate with respect to the variable of the derivative, keeping the other variable fixed. The constant of integration is then actually a function of the other (fixed) variable. Finally, use the second equation to attempt to find this function.

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EXAMPLE 13.14

Consider the vector field F
x y = 2x cos
2yi − 2x2 sin
2y + 4y2 j We want a real-valued function  such that

 = 2x cos
2y

x

and

 = −2x2 sin
2y − 4y2 

y

Choose one of these equations. If we pick the first, then integrate with respect to x, holding y fixed:

x y = 2x cos
2ydx = x2 cos
2y + g
y The “constant” of integration is allowed to involve y because we are reversing a partial derivative, and for any function of y,

2

x cos
2y + g
y = 2x cos
2y

x as we require. We now have 
x y to within some function g
y. From the second equation we need



= −2x2 sin
2y − 4y2 = x2 cos
2y + g
y

y

y Then −2x2 sin
2y − 4y2 = −2x2 sin
2y + g

y so g

y = −4y2  Choose g
y = −4y3 /3 to obtain the potential function 4 
x y = x2 cos
2y − y3  3 It is easy to check that F = &. Is every vector field in the plane conservative? As the following example shows, the answer is no.

EXAMPLE 13.15

Let F
x y =
2xy2 + yi +
2x2 y + ex yj. If this vector field were conservative, there would be a potential  such that



 = 2xy2 + y = 2x2 y + ex y

x

y Integrate the first equation with respect to x to get

x y =
2xy2 + ydx = x2 y2 + xy + f
y

13.3 Independence of Path and Potential Theory in the Plane

539

From the second equation,



= 2x2 y + ex y =
x2 y2 + xy + f
y = 2x2 y + x + f

y

y

y But this would imply that f

y = ex y − x and we cannot find a function of y alone satisfying this equation. Therefore F has no potential. Because not every vector field is conservative, we need some test to determine whether or not a given vector field is conservative. The following theorem provides such a test. THEOREM 13.5

Test for a Convervative Field

Let f and g be continuous in a region D bounded by a rectangle having its sides parallel to the axes. Then F
x y = f
x yi + g
x yj is conservative on D if and only if, for all
x y in D,

g

f = 

x

y

(13.4)

Sometimes the conditions of the theorem hold throughout the plane, and in this event the vector field is conservative for all
x y when equation (13.4) is satisfied.

EXAMPLE 13.16

Consider again F
x y =
2xy2 + yi +
2x2 y + ex yj, from Example 13.15. Compute

f = 4xy + 1

y

and

g = 4xy + ex y

x

and these are unequal on any rectangular region of the plane. This vector field is not conservative. We showed in Example 13.15 that no potential function can exist for this field.

13.3.1

A More Critical Look at Theorem 13.5

The condition (13.4) derived in Theorem 13.5 can be written

g f − = 0

x y But the combination

g f −

x y also occurs in Green’s theorem. This must be more than coincidence. In this section we will explore connections between independence of path, Green’s theorem, condition (13.4), and existence of a potential function. The following example is instructive. Let D consist of all points in the plane except the origin. Thus, D is the plane with the origin punched out. Let −y x F
x y = 2 i+ 2 j = f
x yi + g
x yj 2 x +y x + y2

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Then F is defined on D, and f and g are continuous with continuous partial derivatives on D. Further, we saw in Example 13.12 that

f g − =0

y x

for
x y in D

 Now evaluate C f
x ydx + g
x ydy over two paths from
1 0 to
−1 0. First, let C be the top half of the unit circle, given by x = cos
 y = sin
 for 0 ≤  ≤  Then


=




C

f
x ydx + g
x ydy 

0


− sin

− sin
 + cos
 cos
d =




 0

d = 

Next, let K be the path from
1 0 to
−1 0 along the bottom half of the unit circle, given by x = cos
 y = − sin
 for 0 ≤  ≤  Then


=




K

f
x ydx + g
x ydy

 0

sin

− sin
 + cos

− cos
d = −




 0

d = −

 This means that C f
x ydx + g
x ydy is not independent of path in D. The path chosen between two given points makes a difference. This also means that F is not conservative over D. There is no potential function for F (by Theorem 13.4, if there were a potential function, then the line integral would have to be independent of path). This example suggests that there is something about the conditions specified on the set D in Theorem 13.5 that make a difference. The rectangular set in the theorem, where condition (13.4) is necessary and sufficient for existence of a potential, must have some property or properties that the set in this example lacks. We will explore this line of thought. Let D be a set of points in the plane. We call D a domain if it satisfies two conditions: 1. If P0 is any point of D, there is a circle about P0 such that every point enclosed by this circle is also in D. 2. Between any two points of D, there is a path lying entirely in D. For example, the right quarter plane S consisting of points
x y with x ≥ 0 and y ≥ 0 enjoys property (2), but not (1). There is no circle that can be drawn about a point
x 0 with x ≥ 0, that contains only points with nonnegative coordinates (Figure 13.19). Similarly, any circle drawn about a point
0 y in S must contains points outside of S. The shaded set M of points in Figure 13.20 does not satisfy condition (2). Any path C connecting the indicated points P and Q must at some time go outside of M. Figure 13.21 shows the set A of points between the circles of radius 1 and 3 about the origin. Thus,
x y is in A exactly when 1 < x2 + y2 < 9

13.3 Independence of Path and Potential Theory in the Plane

541

y y (0, y) y

Q A

x

x

C P x

(x, 0) FIGURE 13.19 Right quarter plane x ≥ 0, y ≥ 0.

FIGURE 13.21 The region between two concentric circles.

FIGURE 13.20

This set satisfies conditions (1) and (2), and so is a domain. The boundary circles are drawn as dashed curves to emphasize that points on these curves are not in A. The conditions defining a domain are enough for the first theorem.

THEOREM 13.6

 Let F be a vector field that is continuous on a domain D. Then, C F · dR is independent of path on D if and only if F is conservative.  Proof We know that, if F is conservative, then C F · dR is independent of path on D. It is the converse that uses the condition that D is a domain. Conversely, suppose C F · dR is independent of path on D. We will produce a potential function. Choose any point P0 
x0  y0  in D. If P 
x y is any point of D, define

x y = F · dR C

in which C is any path in D from P0 to P. There is such a path because D is a domain. Further, because this line integral is independent of path, 
x y depends only on
x y and P0 and not on the curve chosen between them. Thus  is a function. Because F is continuous on D,  is also continuous on D. Now let F
x y = f
x yi + g
x yj and select any point
a b in D. We will show that


a b = f
a b and

x


a b = g
a b

y

For the first of these equations, recall that

 
a + x b − 
a b
a b = lim  →0

x x Because D is a domain, there is a circle about
a b enclosing only points of D. Let r be the radius of such a circle and restrict x so that 0 < x < r. Let C1 be any path in D from P0 to
a b and C2 the horizontal line segment from
a b to
a + x b, as shown in Figure 13.22. Let C be the path from P0 to
a + x b consisting of C1 and then C2 . Now



a + x b − 
a b = F · dR − F · dR = F · dR C

C1

C2

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CHAPTER 13

Vector Integral Calculus y

(a, b) C2

y

(a  x, b) x  0

(a  x, b) x 0

C1

(x 0 , y0 )

(x 0 , y0 )

x

FIGURE 13.22

(a, b)

x

FIGURE 13.23

Parametrize C2 by x = a + t x, y = b for 0 ≤ t ≤ 1. Then 
a + x b − 
a b

= F · dR = f
x ydx + g
x ydy = Then




C2

1 0

C2

f
a + t x b
xdt


a + x b − 
a b
1 = f
a + t x bdt x 0

By the mean value theorem for integrals, there is a number  between 0 and 1, inclusive, such that
1 f
a + t x bdt = f
a +  x b 0

Therefore 
a + x b − 
a b = f
a +  x b x As x → 0, f
a +  x b → f
a b by continuity of f , proving that lim

x→0+


a + x b − 
a b = f
a b x

By a similar argument, using the path of Figure 13.23, we can show that lim

x→0−


a + x b − 
a b = f
a b x

Therefore


a b = f
a b

x To prove that
/ y
a b = g
a b, the reasoning is similar except now use the paths of Figures 13.24 and 13.25. This completes the proof of the theorem.

y

(a, b)

(a, b  y) y  0

y (a, b  y), y 0

(x 0 , y0 ) FIGURE 13.24

x

(x 0 , y0 ) FIGURE 13.25

x

13.3 Independence of Path and Potential Theory in the Plane

543

y C x

FIGURE 13.26 The set of points between two concentric circles is not simply connected.

We have seen that the condition (13.4) is necessary and sufficient for F to be conservative within a rectangular region (which is a domain). Although this result is strong enough for many purposes, it is possible to extend it to regions that are not rectangular in shape if another condition is added to the region. A domain D is called simply connected if every simple closed path in D encloses only points of D. A simply connected domain is one that has no “holes” in it, because a simple closed path about the hole will enclose points not in the domain. If D is the plane with the origin removed, then D is not simply connected, because the unit circle about the origin encloses a point not in D. We have seen in Example 13.12 that condition (13.4) may be satisfied in this domain by a vector field having no potential function on D. Similarly, the region between two concentric circles is not simply connected, because a closed curve in this region may wrap around the inner circle, hence enclose points not in the region (Figure 13.26). We will now show that simple connectivity is just what is needed to ensure that condition (13.4) is equivalent to existence of a potential function. The key is that simple connectivity allows the use of Green’s theorem. THEOREM 13.7

Let F
x y = f
x yi + g
x yj be a vector field and D a simply connected domain. Suppose f , g, f/ y and g/ x are continuous on D. Then, F is conservative on D if and only if

f

g =

y

x for all
x y in D. Proof

Suppose first that F is conservative, with potential . Then f
x y =



x

and g
x y =

 

y

Then

f

2 

2 

g = = =

y

x y y x y for
x y in D.  For the converse, suppose that condition (13.4) holds throughout D. We will prove that F · dR is independent of path in D. By the previous theorem, this will imply that F is C conservative. To this end, let P0 and P1 be any points of D, and let C and K be paths in D from P0 to P1 . Suppose first that these paths have only their end points in common (Figure 13.27 (a)). We can then form a positively oriented simply closed path J from P0 to P0 by moving

544

CHAPTER 13

Vector Integral Calculus y
K

y

P1

K P1

D* P0

C1

P0

C1 x

x

FIGURE 13.27(a)

FIGURE 13.27(b)

Paths C and K from P0 to P1 .

A closed path formed from C and −K.

from P0 to P1 along C, then back to P0 along −K (Figure 13.27 (b)). Let J enclose a region D∗ . Since D is simply connected, every point in D∗ is in D, over which f , g, f/ y and g/ x are continuous. Apply Green’s theorem to write

 g f  − F · dR = dA = 0

x y J D∗ But then

J

F · dR = =

so


C





C

C

F · dR + F · dR −

F · dR =


K





−K

K

F · dR

F · dR = 0

F · dR

If C and K intersect each other between P0 and P1 as in Figure 13.28 then this conclusion can still be drawn by considering closed paths between successive points of intersection. We will not pursue this technical argument.  Once we have independence of path of C F · dR on D, then F is conservative, and the theorem is proved. In sum: conservative vector field =⇒ independence of path of


C

F · dR on a set D

independence of path on a domain ⇐⇒ conservative vector field and conservative on a simply connected domain ⇐⇒

y D4

P0

D1

D2

P1

D3

x FIGURE 13.28

f

g = 

x

y

13.4 Surfaces in 3-Space and Surface Integrals

SECTION 13.3

545

PROBLEMS 11. F = 2xyi +
x2 − 1/yj
1 3
2 2 (the path cannot cross the x axis)

In each of Problems 1 through 8, determine whether F is conservative in the given region D. If it is, find a potential function. If D is not defined, it is understood to be the entire plane.

12. F = i +
6y + sin
yj
0 0
1 3 13. F =
3x2 y2 − 6y3 i +
2x3 y − 18xy2 j
0 0
1 1 y 14. F = i + ln
xj
1 1
2 2 (the path must lie in the x right half-plane x > 0)

1. F = y3 i +
3xy2 − 4j 2. F =
6y + yexy i +
6x + xexy j 3. F = 16xi +
2 − y2 j 4. F = 2xy cos
x2 i + sin
x2 j     2x 2y i + j, D the plane with the 5. F = x2 + y 2 x2 + y 2 origin removed

15. F =
−8ey + ex i − 8xey j
−1 −1
3 1   3 16. F = 4xy + 2 i + 2x2 j
1 2
3 3 (the path must x lie in the half-plane x > 0)

6. F = sinh
x + y
i + j

17. Prove the law of conservation of energy: the sum of the kinetic and potential energies of an object acted on by a conservative force field is a constant. Hint: The kinetic energy is
m/2 R

t2 , where m is the mass and R
t the position vector of the particle. The potential energy is −
x y , where  is a potential function for the force. Show that the derivative of the sum of the kinetic and potential energies is zero along any path of motion.

7. F = 2 cos
2xey i + ey sin
2x − yj 8. F =
3x2 y − sin
x + 1i +
x3 + ey j

/ In each of Problems 9 through 16, evaluate C F · dR for C any path from the first given point to the second. 9. F = 3x2
y2 − 4yi +
2x3 y − 4x3 j
−1 1
2 3 10. F = ex cos
yi − ex sin
yj
0 0
2 /4

13.4

Surfaces in 3-Space and Surface Integrals Analogous to the integral of a function over a curve, we would like to develop an integral of a function over a surface. This will require some background on surfaces. A curve is often given by specifying coordinate functions, each of which is a function of a single variable or parameter. For this reason we think of a curve as a one-dimensional object, although the graph may be in 2-space or 3-space. A surface may be defined by giving coordinate functions which depend on two independent variables, say x = x
u v

y = y
u v

z = z
u v

with
u v varying over some set in the u v plane. The locus of such points may form a two-dimensional object in the plane or in R3 .

EXAMPLE 13.17

Suppose a surface is given by the coordinate functions x = au cos
v

y = bu sin
v

z = u

with u and v any real numbers and a and b nonzero constants. In this case it is easy to write z in terms of x and y, a tactic that is sometimes useful in visualizing the surface. Notice that  x 2  y 2 + = cos2
v + sin2
v = 1 au bu

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so x2 y 2 + = u2 = z2  a2 b 2 In the plane y = 0 (the x z plane), z = ±x/a, which are straight lines of slope ±1/a through the origin. In the plane x = 0 (the y z plane), z = ±y/b, and these are straight lines of slope ±1/b through the origin. The surface intersects a plane z = c = constant = 0, in an ellipse x2 y 2 + = c2  a2 b 2

z = c

This surface is called an elliptical cone because it has elliptical cross sections parallel to the x y plane.

EXAMPLE 13.18

Consider the surface having coordinate functions x = u cos
v

y = u sin
v

1 z = u2 sin
2v 2

in which u and v can be any real numbers. Now 1 z = u2 sin
2v = u2 sin
v cos
v 2 = u cos
v u sin
v = xy This surface intersects any plane z = c = constant = 0 in the hyperbola xy = c, z = c. However, the surface intersects a plane y = ±x in a parabola z = ±x2 . For this reason this surface is called a hyperbolic paraboloid. Sometimes x and y are used as parameters, and the surface is defined by giving z as a function of x and y, say z = S
x y. Now the graph of the surface is the locus of points
x y S
x y, as
x y varies over some set of points in the x y plane.

EXAMPLE 13.19

Consider z = write



4 − x2 − y2 for x2 + y2 ≤ 4. By squaring both sides of this equation, we can x2 + y2 + z2 = 4

This appears to be the equation of a sphere of radius 2 about the origin. However, in the original formulation with z given by the radical, we have z ≥ 0, so in fact we have not the sphere, but the hemisphere (upper half of the sphere) of radius 2 about the origin.

EXAMPLE 13.20

 The equation z = x2 + y2 for x2 + y2 ≤ 8 determines a cone having circular cross sections √ parallel to the x y plane. The “top” of the cone is the circle x2 + y2 = 8 in the plane z = 8

13.4 Surfaces in 3-Space and Surface Integrals

547

EXAMPLE 13.21

The equation z = x2 +y2 defines a parabolic bowl, extending to infinity in the positive z-direction because there is no restriction on x or y. These surfaces are easy to visualize and sketch by hand. If the defining function is more complicated then we usually depend on a software package to sketch all or part of the surface. Examples are given in Figures 13.29 through 13.33.

y

2

z

2

1

x

1 0

-1

z

-1

-2

-2 -2

FIGURE 13.29

z = cos
x2 − y2 .

FIGURE 13.30

x 2

y 2 1 -1-2 1 -1 0

6 sin
x − y z=  . 1 + x2 + y 2

-2 z -2

2y

1 0 1 2x -1

-1 -2

FIGURE 13.31

z=

4 cos
x2 + y2  . 1 + x2 + y 2

2x

-1 z -1 -2

1

0 1 2 y

FIGURE 13.32 z = x2 cos
x2 − y2 .

548

CHAPTER 13

Vector Integral Calculus

z

z -2

-2 -1

-1 0 1

y

1

Tv 0

2x

P0

2

∑u 0 y

x FIGURE 13.33

∑v 0

Tu 0

z = cos
xy log
4 + y.

∑

FIGURE 13.34 Tangents to curves  on at P0 determine the tangent plane to the surface there.

Just as we can write a position vector to a curve, we can write a position vector R
u v = x
u vi + y
u vj + z
u vk for a surface. For any u and v in the parameter domain, R
u v can be thought of as an arrow from the origin to the point
x
u v y
u v z
u v on the surface. A surface is simple if R
u1  v1  = R
u2  v2  can occur only if u1 = u2 and v1 = v2 . A simple surface is one that does not fold over and return to the same point for different values of the parameter pairs.

13.4.1 Normal Vector to a Surface

 Let be a surface with coordinate functions x = x
u v, y = y
u v, z = z
u v. Assume that these functions are continuous with  continuous first partial derivatives. Let P0 
x
u0  v0  y
u0  v0  z
u0  v0  be a point on  . If we fix v = v0 , we can define the curve v0 on the surface, having coordinate functions x = x
u v0 

y = y
u v0 

z = z
u v0 

(See Figure 13.34.) This is a curve because its coordinate functions are functions of the single variable u. The tangent vector to v0 at P0 is Tv0 =

x

y

z
u0  v0 i +
u0  v0 j +
u0  v0 k

u

u

u

Similarly, if we fix u = u0 and use v as parameter, we obtain the curve (also shown in Figure 13.34). This curve has coordinate functions x = x
u0  v

y = y
u0  v

z = z
u0  v

 u0

on the surface

13.4 Surfaces in 3-Space and Surface Integrals The tangent vector to

 u0

549

at P0 is

Tu0 =

x

y

z
u0  v0 i +
u0  v0 j +
u0  v0 k

v

v

v

Assumingthat neither of these tangent vectors is the zero vector, they both lie in the tangent plane to at P0 . Their cross product is therefore  normal (orthogonal) to this tangent plane, and is the vector we define to be the normal to at P0 :

x

y

z
u0  v0 i +
u0  v0 j +
u0  v0 k

u

u

u

x

y

z
u  v i +
u0  v0 j +
u0  v0 k ×

v 0 0

v

v     i j k     x

y

z  =  u
u0  v0  u
u0  v0  u
u0  v0      x 
u0  v0  y
u0  v0  z
u0  v0  

v

v

v      

y z z y

z x x z

x y y x = − − − i+ j+ k

u v u v

u v u v

u v u v

N
P0  =

(13.5)

in which all the partial derivatives are evaluated at
u0  v0 . An expression that is easier to remember is obtained by introducing Jacobian notation. Define the Jacobian determinant (named for the German mathematician Karl Jacobi) of functions f and g to be the 2 × 2 determinant 


f g  =


u v  In this notation, the normal vector to

N
P0  =

f

u

f

v

g

u

g

v



  f g g f  −  =  u v u v

at P0 is


z x


x y


y z i+ j+ k


u v


u v


u v

with all the partial derivatives evaluated at
u0  v0 . This notation helps in remembering the normal vector because of the cyclic pattern in the Jacobian symbols. Write x y z in this order. For the first component of N
P0 , omit the first letter, x, to obtain
y z/
u v. For the second component, omit y, but maintain the same cyclic direction, moving left to right through x y z. This means we start with z, the next letter after y, then back to x, obtaining


z x/
u v. For the third component, omit z, leaving x y and the Jacobian


x y/
u v.  at P0 . Of course, any nonzero real multiple of N
P0  is also a normal to

550

CHAPTER 13

Vector Integral Calculus

EXAMPLE 13.22

Consider again the elliptical cone x = au cos
v

y = bu sin
v z = u √ Suppose we want the normal vector at P0 
a 3/4 b/4 1/2, obtained when u = u0 = 1/2, v = v0 = /6. Compute the Jacobians:

y z z y


y z − =


u v
1/2/6

u v u v
1/2/6 √ = b sin
v
0 − bu cos
v
1/2/6 = − 3b/4

z x x z


z x − =


u v
1/2/6

u v u v
1/2/6 = −au sin
v − a cos
v
0
1/2/6 = −a/4 and


x y


u v



x y y x − =

u v u v
1/2/6


1/2/6

= a cos
vbu cos
v − b sin
v
−au sin
v
1/2/6 = ab/2 The normal vector at P0 is √ b a ab N
P0  = − 3 i − j + k 4 4 2 Consider the special case that the surface  is given explicitly as z = S
x y. We may think of u = x and v = y as the parameters for and write the coordinate functions as x = x

y = y

z = S
x y

Since x/ x = 1 = y/ y and x/ y = y/ x = 0, we have  

S


y z  0 1 


y z = = 

S  = − 


u v
x y  S

x

x

y 

S 

S


z x  S

x

y  =−  =


x y  1 0 

y and




x y  1 =


x y  0

 0   = 1 1 

The normal at a point P0 
x0  y0  S
x0  y0  in this case is N
P0  = −

S

S
x  y i −
x0  y0 j + k

x 0 0

y

=−

z

z
x0  y0 i −
x0  y0 j + k

x

y

We can also denote this vector as N
x0  y0 .

(13.6)

13.4 Surfaces in 3-Space and Surface Integrals

EXAMPLE 13.23

Consider the cone given by z = S
x y = x

S = 2

x x + y2



551

x2 + y2 . Then

y

S =  2

y x + y2 √ except for x = y = 0. Take, for example, the point
3 1 10. The normal vector at this point is √ 3 1 N
3 1 10 = − √ i − √ j + k 10 10 and

This normal vector is shown in Figure 13.35, and it points into the cone. In some contexts we want to know whether a normal vector is an inner normal (such as this one) or an outer normal (pointing out of the region bounded by the surface). If we wanted an outer normal at this point, we could take 1 3 −N
P0  = √ i + √ j − k 10 10 This cone does not have a normal vector at the origin, which is a “sharp point” of the surface. There is no tangent plane at the origin.

z

(3, 1,
10)

N

y x FIGURE 13.35  Normal to

the cone z = √
3 1 10.

x2 + y2 at

The normal vector (13.6)  could also have been derived using the gradient vector. If given by z = S
x y, then is a level surface of the function



is


x y z = z − S
x y The gradient of this function is a normal vector, so compute & =





 i+ j+ k

x

y

z

=−

13.4.2

S

S i − j + k = N
P

x

y

The Tangent Plane to a Surface

 If a surface has a normal  vector N at a point P0 , then we can use N to determine the equation of the tangent plane to at P0 . Let
x y z be any point on the tangent plane. Then the vector
x − x0 i +
y − y0 j +
z − z0 k is in the tangent plane, hence is orthogonal to N. Then N ·
x − x0 i +
y − y0 j +
z − z0 k = 0

552

CHAPTER 13

Vector Integral Calculus

More explicitly,




y z


z x


x y
x − x0  +
y − y0  +
z − z0  = 0


u v
u0 v0 


u v
u0 v0 


u v
u0 v0   This is the equation of the tangent plane to at P0 .

EXAMPLE 13.24

Consider again the elliptical cone given by x = au cos
v

y = bu sin
v z = u √ √ We found in Example 13.22 that the normal vector at P 
a 3/4 b/4 1/2 is N = − 3
b/4i− 0 
a/4j +
ab/2k. The tangent plane to at this point has equation  √    √ b b ab 1 a 3 a x− +
z −  = 0 − 3 x− − 4 4 4 4 2 2  In the special case that is given by z = S
x y, then the normal vector at P0 is N = − S/ x
x0  y0 i − S/ y
x0  y0 j + k, so the equation of the tangent plane becomes −

S

S
x0  y0 
x − x0  −
x0  y0 
y − y0  +
z − z0  = 0

x

y

This equation is usually written z − z0 =

S

S
x  y 
x − x0  +
x0  y0 
y − y0 

x 0 0

y

13.4.3 Smooth and Piecewise Smooth Surfaces Recall that a curve is smooth if it has a continuous tangent. Similarly, a surface is smooth if it has a continuous normal vector. A surface is piecewise smooth if it consists of a finite number of smooth surfaces. For example, a sphere is smooth, and the surface of a cube is piecewise smooth. A cube consists of six square pieces, which are smooth, but does not have a normal vector along any of its edges.  In calculus, it is shown that the area of a smooth surface given by z = S
x y is the integral 0  2  2 



S

S area of = 1+ + dA (13.7)

x

y D where D is the set of points in the x y− plane for which S is defined. This may also be written 0   2  2 



z

z area of = 1+ + dxdy

x

y D Equation (13.7) is the integral of the length of the normal vector (13.6): 

area of = N
x y dxdy D

This is analogous to the formula for the length of a curve as the integral of the length of the tangent vector.

13.4 Surfaces in 3-Space and Surface Integrals

553

 More generally, if is given by coordinate functions x = x
u v, y = y
u v and z = z
u v, with
u v varying over some set D in the u v plane, then 

(13.8) area of = N
u v dudv D

the integral of the length of the normal vector, which is given by equation (13.5).

EXAMPLE 13.25

We will illustrate these formulasfor surface area for a simple case in which we know the area from elementary geometry. Let be the upper hemisphere  of radius 3 about the origin.  We can write as the graph of z = S
x y = 9 − x2 − y2 , with x2 + y2 ≤ 9. D consists of all points on or inside the circle of radius 3 about the origin in the x y-plane. We can use equation (13.7). Compute x x

z = − =− 2 2

x z 9−x −y and, by symmetry,

z y =− 

y z Then area of



= =



D



D

0 1+ 0

 2   2 x y + dxdy z z



z2 + x 2 + y 2 3 dxdy = dxdy  2 z D 9 − x2 − y 2

This is an improper double integral which we can evaluate easily by converting it to polar coordinates. Let x = r cos
, y = r sin
. Since D is the disk of radius 3 about the origin, 0 ≤ r ≤ 3 and 0 ≤  ≤ 2. Then


2
3 3 3 dxdy = rdrd  √ 2 2 0 0 D 9 − r2 9−x −y
3 3

r = 6 dr = 6 −
9 − r 2 1/2 0 √ 2 0 9−r

1/2  = 18 = 6 9 This is the area of a hemisphere of radius 3. We are now prepared to define the integral of a function over a surface.

13.4.4

Surface Integrals

The notion of the integral of a function over a surface is modeled after the line integral, with respect to arc length, of a function over a curve. Recall that, if a smooth curve C is given by x = x
t, y = y
t, z = z
t for a ≤ t ≤ b, then the arc length along C is
t s
t = x

2 + y

2 + z

2 d a

554

CHAPTER 13

Vector Integral Calculus

Then ds =



x

t2 + y

t2 + z

t2 dt

and the line integral of a function f along C, with respect to arc length, is
C

f
x y zds =




b a

 f
x
t y
t z
t x

t2 + y

t2 + z

t2 dt

We want to lift these ideas up one dimension to integrate over a surface instead of a curve. Now we have coordinate functions that are functions of two independent variables, say u and v, with
u v varying over some given set D in the u v plane. This means that b  will be replaced by D . The differential element of arc length, ds, which is used in the a line integral, will be replaced by the differential element d of surface area on the surface. By equation (13.8), d = N
u  v dudv, in which N
u v is the normal vector at the point
x
u v y
u v z
u v on .

DEFINITION 13.6

Surface Integral



Let be a smooth surface having coordinate functions x = x
u v, y = y
u v, z = z
u v for
u v in some  set D of the u v plane. Let f be continuous on . Then the surface  integral of f over is denoted  f
x y zd, and is defined by



f
x y zd = f
x
u v y
u v z
u v N
u v dudv  D

   If is a piecewise smooth surface having smooth components 1      n , with each component either disjoint from the others, or intersecting another component in a set of zero area (for example, along a curve), then

 f
x y zd

=





f
x y zd + · · · +





1



f
x y zd

n

For example, we would integrate over the surface of a cube by summing the surface integrals over the six faces. Two such faces either do not intersect, or intersect each other along a line segment having zero area. The intersection condition is to prevent the selection of surface components that overlap each other in significant ways. This is analogous to a piecewise smooth curve C formed as the join of smooth curves C1      Cn . When we do this, we assume that two of these component curves either do not intersect, or intersect just at an end point, not along an arc  of both curves. If is described by z = S
x y, then



 f
x y zd =

0



D

f
x y S
x y 1 +



S

x



2 +

S

y

2 dxdy

We will look at some examples of evaluation of surface integrals, then consider uses of surface integrals.

13.4 Surfaces in 3-Space and Surface Integrals

555

z (0, 0, 4) (0, 1, 3) Part of the plane x y z  4

(2, 0, 2)

(0, 1)

(2, 1, 1) (2, 0)

y

(0, 4, 0) (2, 1)

(4, 0, 0) x FIGURE 13.36

EXAMPLE 13.26

  Evaluate  zd if is the part of the plane x + y + z = 4 lying above the rectangle 0 ≤ x ≤ 2, 0 ≤ y ≤ 1. The surface is shown in Figure 13.36. D consists of all
x y with 0 ≤ x ≤ 2 and 0 ≤ y ≤ 1. With z = S
x y = 4 − x − y we have



zd = z 1 +
−12 +
−12 dxdy  D

√
2
1 = 3
4 − x − ydydx 0

0

First compute
0

1

1
4 − x − ydy =
4 − xy − y2 2

1 0

1 7 = 4 − x − = − x 2 2 Then





 zd =

  √ √
2 7 − x dx = 5 3 3 2 0

EXAMPLE 13.27

Recall the hyperbolic paraboloid of Example 13.18 given by 1 y = u sin
v z = u2 sin
2v 2   xyzd over the part of this surface corresponding to We will compute the surface integral 1 ≤ u ≤ 2, 0 ≤ v ≤ . First we need the normal vector. The components of N
u v are the Jacobians:  


y z  sin
v u cos
v  =


u v  u sin
2v u2 cos
2v  x = u cos
v

= u2 sin
v cos
2v − cos
v sin
2v 

556

CHAPTER 13

Vector Integral Calculus 


z x  u sin
2v =


u v  cos
v

 u2 cos
2v  −u sin
v 

= −u2 sin
v sin
2v + cos
v cos
2v  and




x y  cos
v −u sin
v =


u v  sin
v u cos
v

   = u 

Then N
u v2 = u4 sin
v cos
2v − cos
v sin
2v2 + u4 sin
v sin
2v + cos
v cos
2v2 + u2 = u2
1 + u2  so

 N
u v = u 1 + u2 

The surface integral is





1 2 xyzd =

u cos
v u sin
v sin
2v u 1 + u2 dA u  2 D

2  = cos
v sin
v sin
2vdv u5 1 + u2 du 0

=

 4



100 √ 21

 11 √ 5− 2  105

1

As we expect of an integral,






f
x y z + g
x y zd =
f
x y zd +    g
x y zd and, for any real number ,



 f
x y zd = 



 f
x y zd

The next section is devoted to some applications of surface integrals.

PROBLEMS

SECTION 13.4 In 

each of Problems  f
x y zd.

1

through

10,

evaluate

 1. f
x y z = x, is the part of the plane x +4y +z = 10 in the first octant  2. f
x y z = y2 , is the part of the plane z = x for 0 ≤ x ≤ 2, 0 ≤ y ≤ 4 

3. f
x y z = 1, is the part of the paraboloid z = x2 + y2 lying between the planes z = 2 and z = 7

 4. f
x y z = x + y, is the part of the plane 4x + 8y + 10z = 25 lying above the triangle in the x y− plane having vertices
0 0,
1 0 and
1 1   5. f
x y z = z, is the part of the cone z = x2 + y2 lying in the first octant and between the planes z = 2 and z = 4  6. f
x y z = xyz, is the part of the plane z = x + y with
x y lying in the square with vertices
0 0,
1 0,
0 1 and
1 1

13.5 Applications of Surface Integrals

 9. f
x y z = z, is the part of the plane z = x − y for 0 ≤ x ≤ 1 0 ≤ y ≤ 5  10. f
x y z = xyz, is the part of the cylinder z = 1 + y2 for 0 ≤ x ≤ 1 0 ≤ y ≤ 1

 7. f
x y z = y, is the part of the cylinder z = x2 for 0 ≤ x ≤ 2 0 ≤ y ≤ 3  8. f
x y z = x2 , is the part of the paraboloid z = 4 − x2 − y2 lying above the x y plane

13.5

557

Applications of Surface Integrals 13.5.1 If



Surface Area

is a piecewise smooth surface, then



 1d = . N
u v dudv = area of  D

(This assumes a bounded surface having finite area.) Clearly we do not need the notion of a surface integral to compute the integral and obtain the area of a surface. However, we mention this result because it is in the same spirit as other familiar mensuration formulas:
ds = length of C



C

D

dA = area of D

and, if M is a solid region in 3-space enclosing a volume, then


dV = volume of M. M

13.5.2

Mass and Center of Mass of a Shell

 Imagine a thin shell of negligible thickness taking the shape of a smooth surface . Let 
x y z be the density of the material of the shell at
x y z. Assume that  is continuous. We want to compute the mass of the shell.  Suppose has coordinate functions x = x
u v, y = y
u v, z = z
u v for
u v in D. Form a grid over D in the u v plane by drawing lines (dashed lines in Figure 13.37) parallel to the axes, and retain only those rectangles R1      RN intersecting D. Let the vertical lines be u units apart, and the horizontal lines, v units apart. For
u vvarying over  Rj , we obtain a patch v for
u v in Rj . or surface element j on the surface (Figure 13.38). That is, j
u v =
u Let
uj  vj  be a point in Rj , and approximate the density of the surface element j by the constant j = 
x
uj  vj  y
uj  vj  z
uj  vj   Because  is continuous and j has finite area, we can choose u and v sufficiently small  as closely as we like over  j . that j approximates 
x y z Approximate the mass of j as j times the area of j . Now this area is

 area of j = N
u v dudv Rj

# # ≈ #N
uj  vj # u v so mass of The mass of



 j

# # ≈ j #N
uj  vj # u v

is approximately the sum of the masses of the surface elements:

mass of



≈

N  j=1

# # 
x
uj  vj  y
uj  vj  z
uj  vj  #N
uj  vj # u v

558

CHAPTER 13

Vector Integral Calculus

In 1927, Congress approved construction of a dam to control the Colorado River for the purpose of fostering agriculture in the American southwest and as a source of hydroelectric power, spurring the growth of Las Vegas and southern California. Construction began in 1931. One major problem of construction was the cooling of concrete as it was poured. Engineers estimated that the amount of concrete required would take 100 years to cool. The solution was to pour it in rows and columns of blocks, through which cooled water was pumped in pipes. Hoover Dam was completed in 1935 and is 727 feet high, 1244 feet long, 660 feet thick at its base, and 45 feet thick at the top. It weighs about 5.5 million tons and contains 3,250,000 cubic yards of concrete. On one side of the dam, Lake Mead is over 500 feet deep. Computation of forces and stresses on parts of the dam surface involve a combination of material science, fluid flow, and vector analysis.

v z (x(uj , vj), y(uj , vj), z(uj , vj)) Rj ∆v

(uj , vj)

∑j

∑

u y ∆u FIGURE 13.37

plane.

Grid over D in the u, v

x FIGURE 13.38

Surface element

 j

on



.

13.5 Applications of Surface Integrals

559

 This is a Riemann sum for D 
x
u v y
u v z
u v N
u v dudv. Hence in the limit as u → 0 and v → 0 we obtain 

mass of =  
x y zd The mass of the shell is the surface integral of the density function. This is analogous to the mass of a wire being the line integral of the density function over the wire. The center of mass of the shell is
x y z, where 1

1

x= x
x y zd y = y
x y zd m  m  and z=

1

z
x y zd m 

in which m is the mass of the shell.

EXAMPLE 13.28

 We will find the mass and center of mass of the cone z = x2 + y2 , where x2 + y2 ≤ 4 and the density function is 
x y z = x2 + y2 . First calculate the mass m. We will need

z x

z y = and = 

x z

y z Then m= = =




x

D


x2 + y2  1 +

2
2 0

+ y2 d 0






2

x2 y 2 + dydx z2 z2

√ r 2 2rdrd

0

√ 1 = 2 2 r 4 4

2

√ = 8 2

0

By symmetry of the surface and of the density function, we expect the center of mass to lie on the z axis, so x = y = 0. Finally, 1

2 2 z= √  z
x + y d 8 2 0 x2 y 2 1

 2 = √ x + y2
x2 + y2  1 + 2 + 2 dydx z z 8 2 D

2 2 1 = r
r 2 rdrd 8 0 0

1 5 2 8 1 =
2 r =  8 5 5 0 The center of mass is
0 0 85 .

560

CHAPTER 13

Vector Integral Calculus

13.5.3 Flux of a Vector Field Across a Surface Suppose a fluid moves in some region of 3-space  (for example, through a pipeline), with velocity V
x y z t. Consider an imaginary surface within the fluid, with continuous unit  normal vector N
u v t. The flux of V across is the net volume of fluid, per unit time,  flowing across in the direction of N. We would like to calculate this flux.   In a time interval t the volume equals the  of fluid flowing across a small piece j of V t, where V is the component of V in volume of the cylinder with base j and altitude N  N . This volume (Figure 13.39) is
V the direction of N, evaluated at some point of N tAj , j  is the area of . Because = 1, V = V · N. The volume of fluid flowing across where A N j N j  per unit time is j VN
tAj = VN Aj = V · NAj  t Sum these quantities over all the pieces of the surface and take the limit as the pieces are chosen smaller, as we did for the mass of the shell. We get

 flux of V across in the direction of N =  V · Nd The flux of a vector across a surface is therefore computed as the surface integral of the normal component of the vector to the surface.

V

Volume  VN t Area of  j 

z

N



VN t j y

x FIGURE 13.39

EXAMPLE 13.29

Find the flux of F = xi + yj + zk across the part of the sphere x2 + y2 + z2 = 4 lying between the planes z = 1  and z = 2. The surface is shown in Figure 13.40, along with the normal vector (computed below) at a point. We may think of as defined by z = S
x y , where S is defined by the equation of the sphere and
x y varies over a set D in the x y plane. To determine D, observe that the plane z = 2 hits only at its “north pole”
0 0 2. The plane z = 1 intersects the sphere √ in the circle x2 + y2 = 3 , z = 1. This circle projects onto the x y plane to the circle of radius 3 about the origin. Thus D consists of points
x y satisfying x2 + y2 ≤ 3 (shaded in Figure 13.41). To compute the partial derivatives z/ x and z/ y, we can implicitly differentiate the equation of the sphere to get 2x + 2z

z = 0

x

13.5 Applications of Surface Integrals z

z

x2  y2  z2  4 1  z 2

(0, 0, 2)
12

561

(x i  y j  zk)

(0, 0, 1)

y y

x

x

FIGURE 13.40

x2  y2  3

FIGURE 13.41

so x

z =−

x z and, similarly, y

z =− 

y z A normal vector to the sphere is therefore x y − i − j − k z z Since we need a unit normal in computing flux, we must divide this vector by its length, which is 0 0 x2 y 2 x 2 + y 2 + z2 2 + 2 +1 = =  2 2 z z z z A unit normal is therefore

  z y 1 x N= − i − j − k = −
xi + yj + zk  2 z z 2  This points into the sphere. If we want the flux across from the  outside of the sphere toward the inside, this is the normal to use. If we want the flux across from within the sphere, use −N instead. Now,  1 2 F ·
−N = x + y 2 + z2  2 There

1 flux = 
x2 + y2 + z2 d 2 0 x2 y 2 1

2
x + y2 + z2  1 + 2 + 2 dA = 2 D z z 0 1

2 x 2 + y 2 + z2 =
x + y2 + z2  dA 2 D z2 3/2 1 1

 2 x + y 2 + z2 dA =  2 D 4 − x2 − y 2

1 =4 dA  D 4 − x2 − y 2

CHAPTER 13

562

Vector Integral Calculus

because x2 + y2 + z2 = 4 on

 . Converting to polar coordinates, we have


2


√

3

1 rdrd √ 4 − r2

√3 = 8 −
4 − r 2 1/2 0 = 8

flux = 4

0

0

We will see other applications of surface integrals when we discuss the integral theorems of Gauss and Stokes, and again when we study partial differential equations.

SECTION 13.5

PROBLEMS

1. Find the mass and center of mass of the triangular shell having vertices
1 0 0,
0 3 0 and
0 0 2 if 
x y z = xz + 1. 2. Find the center of mass of the portion of the homogeneous sphere x2 + y2 + z2 = 9 lying above the plane z = 1. (Homogeneous means that the density function is constant).

2 2 2 2 tween the cylinders  x + y = 1 and x + y = 9, if 
x y z = xy/ 1 + 4x2 + 4y2 .

5. Find the mass and center of mass  of the paraboloid z = 6 − x2 − y2 if 
x y z = 1 + 4x2 + 4y2 . 6. Find the center of mass of the part of the homogeneous sphere x2 + y2 + z2 = 1 lying in the first octant.

3. Findthe center of mass of the homogeneous cone z = x2 + y2 for x2 + z2 ≤ 9.

7. Find the flux of F = xi + yj − zk across the part of the plane x + 2y + z = 8 lying in the first octant.

4. Find the center of mass of the part of the paraboloid z = 16 − x2 − y2 lying in the first octant and be-

8. Find the flux of F = xzi − yk across the part of the sphere x2 + y2 + z2 = 4 lying above the plane z = 1.

13.6

Preparation for the Integral Theorems of Gauss and Stokes The fundamental results of vector integral calculus are the theorems of Gauss and Stokes. In this section we will begin with Green’s theorem and explore how natural generalizations lead to these results. With appropriate conditions on the curve and the functions, the conclusion of Green’s theorem is

 g f  − f
x ydx + g
x ydy = dA

x y C D in which D is the region on and enclosed by the simple closed smooth curve C. Define the vector field F
x y = g
x yi − f
x yj Then & ·F =

g f − 

x y

Now parametrize C by arc length, so the coordinate functions are x = x
s, y = y
s for 0 ≤ s ≤ L. The unit tangent vector to C is T
s = x

si + y

sj, and the unit normal vector

13.6 Preparation for the Integral Theorems of Gauss and Stokes

563

is N
s = y

si − x

sj. These are shown in Figure 13.42. This normal points away from the interior D of C, and so is an outer normal. Now F · N = g
x y so

C

dx dy + f
x y  ds ds

- dx dy f
x y + g
x y ds ds ds C = F · Nds

f
x ydx + g
x ydy =

C

We may therefore write the conclusion of Green’s theorem in vector form as

F · Nds = & · FdA C

(13.9)

D

This is a conservation of energy equation. Recall from Section 12.4.1 that the divergence of a vector field at a point is a measure of the flow of the field from that point. Equation (13.9) states that the flux of the vector field outward from D across C (because N is an outer normal) exactly balances the flow of the field from each point in D. The reason for writing Green’s theorem in this form is that it suggests a generalization to  three dimensions. Replace the closed curve C in the plane with a closed surface in 3-space (closed  meaning bounding a volume). Replace the line integral over C with a surface integral over , and allow the vector field F to be a function of three variables. We conjecture that equation (13.9) generalizes to




& · FdV  F · Nd = M



 in which N is a unit normal to pointing away from the solid region M bounded by . We will see that, under suitable conditions on and F, this is the conclusion of Gauss’s divergence theorem. Now begin again with Green’s theorem. We will pursue a different generalization to three dimensions. This time let F
x y z = f
x yi + g
x yj + 0k The reason for adding the third component is to be able to take the curl:    i j k    

g f   & × F =  x y z  = − k  

x y  f g 0  Then
& × F · k =

y

g f − 

x y

z T N

 C

x FIGURE 13.42

x FIGURE 13.43

N

y

CHAPTER 13

564

Vector Integral Calculus

Further, with unit tangent T
s = x

si + y

sj to C, we can write   dx dy F · Tds = f
x yi + g
x yj · i + j ds ds ds = f
x ydx + g
x ydy so the conclusion of Green’s theorem can also be written

F · Tds =
& × F · kdA C

(13.10)

D

Now think of D as a flat surface in the x y plane, with unit normal vector k, and bounded by the closed curve C. To generalize this, allow C to be a curve in 3-space bounding a surface  having unit outer normal vector N, as shown in Figure 13.43. Now equation (13.10) suggests that

F · Tds = 
& × F · Nd C

We will see this equation shortly as the conclusion of Stokes’s theorem.

SECTION 13.6

PROBLEMS

1. Let C be a simple closed path in the x y plane, with interior D. Let 
x y and 
x y be continuous with continuous first and second partial derivatives on C and throughout D. Let & 2 = Prove that



=

-

C

& 2 dA

D

C

2  2  + 2

x2

y

3. Let C be a simple closed path in the x y plane, with interior D. Let  be continuous with continuous first and second partial derivatives on C and at all points of D. Let N
x y be the unit outer normal to C (outer meaning pointing away from D if drawn as an arrow at
x y on C). Prove that

N
x yds = & 2 
x ydA

−

D

(Recall that N
x y is the directional derivative of  in the direction of N).







 dx +  dy − & · &dA

y

x D

2. Under the conditions of Problem 1, show that

  & 2  − & 2  dA D



- 





 = − dx +  − dy 

y

y

x

x C

13.7

The Divergence Theorem of Gauss We have seen that, under certain conditions, the conclusion of Green’s theorem is

F · Nds = & · FdA C

D

13.7 The Divergence Theorem of Gauss

565

Now make the following generalizations from the plane to 3-space: set D in the plane → a 3-dimensional solid M  a closed curve C bounding D → a surface enclosing M  a unit outer normal N to C → a unit outer normal N to a vector field F in the plane→ a vector field F in 3-space /  a line integral C F · Nds → a surface integral  F · Nd   & · FdV a double integral D & · FdA → a triple integral M With these correspondences and some terminology, Green’s theorem suggests a theorem named for the great nineteenth-century German mathematician and scientist Carl Friedrich Gauss.  A surface is closed if it encloses a volume. For example, a sphere is closed, as is a cube, while a hemisphere is not. A surface consisting of the top part of the sphere x2 + y2 + z2 = a2 , together with the disk x2 + y2 ≤ a2 in the x y plane, is closed. If filled with water (through some opening that is then sealed off), it will hold the water.  A normal vector N to is an outer normal if, when represented as an arrow from a point of the surface, it points away from the region enclosed by the surface (Figure 13.44). If N is also a unit vector, then it is a unit outer normal. z

N (Unit outer normal)

M 

y x FIGURE 13.44

THEOREM 13.8

Gauss’s Divergence Theorem

 Let be  a piecewise smooth closed surface. Let M be the set of points on and enclosed by . Let have unit outer normal vector N. Let F be a vector field  whose components are continuous with continuous first and second partial derivatives on and throughout M. Then




& · FdV (13.11)  F · Nd = M

& · F is the divergence of the vector field, hence the name “divergence theorem”. In the spirit of Green’s theorem, Gauss’s theorem relates vector operations over objects of different dimensions. A surface is a two-dimensional object (it has area but no volume), while a solid region in 3-space is three-dimensional. Gauss’s theorem has several kinds of applications. One is to replace one of the integrals in equation (13.11) with the other, in the event that this simplifies an integral evaluation. A second is to suggest interpretations of vector operations. A third is to serve as a tool in deriving physical laws. Finally, we will use the theorem in developing relationships to be used in solving partial differential equations. Before looking at uses of the theorem, here are two purely computational examples to provide some feeling for equation (13.11).

566

CHAPTER 13

Vector Integral Calculus

EXAMPLE 13.30

  Let  be the piecewise smooth closed surface consisting of the surface 1 of the cone  z = x2 + y2 for x2 + y2 ≤ 1, together with the flat cap 2 consisting of the disk x2 + y2 ≤ 1 in the plane z = 1. This surface is shown in Figure 13.45. Let F
x y z = xi + yj + zk. We will calculate both sides of equation  (13.11). The unit outer normal to 1 is   y 1 x i+ j−k  N1 = √ z 2 z Then

because on

  y 1 x F · N1 =
xi + yj + zk · √ i+ j−k z 2 z   2 2 y 1 x + −z = 0 =√ z z 2 

1,

z2 = x2 + y2 . (One can also see geometrically that F is orthogonal to N1 .) Then

 F · N1 d = 0

The unit outer normal to

1

 2

is N2 = k, so F · N2 = z

Since z = 1 on



2,

then



F · N2 d =





2

2

= area of Therefore





 F · Nd =









zd =

F · N1 d +





= 

2





1

F · N2 d = 

2

Now compute the triple integral. The divergence of F is & ·F =



x + y + z = 3

x

y

z

z

(0, 0, 1) 2 1

x2  y2  z 2 y

x FIGURE 13.45

d

2

13.7 The Divergence Theorem of Gauss so




M

567






& · FdV =

3dV M

= 3 volume of the cone of height 1, radius 1 1 = 3  =  3

EXAMPLE 13.31

Let



be the piecewise smooth surface of the cube having vertices
0 0 0
1 0 0
0 1 0
0 0 1
1 1 0
0 1 1
1 0 1
1 1 1

Let F
x y z = x2 i + y2 j + z2 k. We would like to compute the flux of this vector field across the faces of the cube.  The flux is  F · Nd. This integral can certainly be calculated directly, but it requires  performing the integration over each of the six smooth faces of . It is easier to use the triple integral from Gauss’s theorem. Compute the divergence & · F = 2x + 2y + 2z and then flux = = = = =





 F · Nd




M

& · FdV = 2


1
1
0


1
0




0

0

1

0 1

0

1




M


x + y + zdV


2x + 2y + 2zdzdydx


2x + 2y + 1dydx


2x + 2dx = 3

Now we will move to more substantial uses of the theorem.

13.7.1

Archimedes’s Principle

Archimedes’s Principle states that the buoyant force a fluid exerts on a solid object immersed in it, is equal to the weight of the fluid displaced. A bar of soap floats or sinks in a full bathtub for the same reason a battleship floats or sinks in the ocean. The issue rests with the weight of the fluid displaced by the object. We will derive this principle.  Consider a solid object M bounded by a piecewise smooth surface . Let  be the constant density of the fluid. Draw a coordinate system as in Figure 13.46, with M below the surface of the fluid. Using  the fact that pressure is the product of depth and density, the pressure p
x y z at a point on is given by p
x y z = −z. The negative sign is used because z is negative in the downward direction and  we want pressure to be positive. the Now consider a piece j of the surface, also shown in Figure 13.46. The force of  pressure on this surface element has magnitude approximately −z times the area Aj of j .

568

CHAPTER 13

Vector Integral Calculus z y k x

N j

M FIGURE 13.46

Pressure force on  j ≈ zN · Aj . Vertical component = zN · kAj .

  If N is the unit outer normal to j , then the force caused by the pressure on j is approximately zNA  j . The vertical component of this force is the magnitude of the buoyant force acting upward on j . This vertical component is zN · kAj . Sum these vertical components over the entire surface to obtain approximately the net buoyant force on the object, then take the limit as the surface elements are chosen smaller (areas tending to zero). We obtain in this limit that 

net buoyant force on =  zN · kd Write this integral as into a triple integral:



 zk · Nd

and apply Gauss’s theorem to convert the surface integral

net buoyant force on But & ·
zk = , so net buoyant force on



=



=




M




M

& ·
zkdV

dV =  volume of M 

But this is exactly the weight of the fluid displaced, establishing Archimedes’s Principle.

13.7.2 The Heat Equation We will derive a partial differential equation that models heat conduction. Suppose some medium (for example, a metal bar, the air in a room or water in a pool) has density 
x y z, specific heat 
x y z, and coefficient of thermal conductivity K
x y z. Let u
x y z t be the temperature of the medium at time t and point
x y z. We want to derive an equation for u. We will employ a device used  frequently in deriving mathematical models. Consider an imaginary smooth closed surface  within the medium, bounding a solid region M. The amount of heat energy leaving M across in a time interval t is  


K&u · Nd t This is the flux of the vector (K times the gradient of u) across this surface, multiplied by the length of the time interval. But, the change in temperature at
x y z in M in this time interval is approximately
u/ t t, so the resulting heat loss in M is 




u  dV t

t M

13.7 The Divergence Theorem of Gauss

569

Assuming that there are no heat sources or losses within M (for example, chemical reactions or radioactivity),  the change in heat energy in M over this time interval must equal the heat exchange across . Then 

 




u
K&u · Nd t =  dV t 

t M Therefore




K&u · Nd =




 M

u dV

t

Apply Gauss’s theorem to the surface integral to obtain







u & ·
K&udV =  dV

t M M The role of Gauss’s theorem here is to convert the surface integral to a triple integral, thus obtaining an equation with the same kind of integral on both sides. This allows us to combine terms and write the last equation as 


 u  − & ·
K&u dV = 0

t M  Now keep in mind a crucial point— is any smooth closed surface within the medium. Assume that the integrand in the last equation is continuous. If this integrand were nonzero at any point P0 of the medium, then it would be positive or negative at P0 , say positive. By continuity of this integrand, there would be a sphere S, centered at P0 , of small enough radius that the integrand would be strictly positive on and within S. But then we would have 


 u  − & ·
K&u dV > 0

t M  in which M is the solid ball bounded by S. By choosing = S, this is a contradiction. We conclude that 

u − & ·
K&u = 0

t

at all points in the medium, for all times. This gives us the partial differential equation 

u = & ·
K&u

t

for the temperature function at any point and time. This equation is called the heat equation. We can expand  

u

u

u & ·
K&u = & · K i + K j + K k

x

y

z

u

u

u
K  +
K  +
K 

x x

y y

z z  2 

u 2 u 2 u

K u K u K u + + +K = + +

x x y y

z z

x2 y2 z2

=

= &K · &u + K& 2 u

570

CHAPTER 13

Vector Integral Calculus

in which & 2u =

2 u 2 u 2 u + +

x2 y2 z2

is called the Laplacian of u. (& 2 is read “del squared”). Now the heat equation can be written 

u = &K · &u + K& 2 u

t

If K is constant, then its gradient is the zero vector and this equation simplifies to K 2

u = & u

t  In the case of one space dimension (for example, if u
x t is the temperature distribution in a thin bar lying along a segment of the x axis), this is

2 u

u =k 2

t

x with k = K/p.K The steady-state case occurs when u does not change with time. In this case u/ t = 0, and the last equation becomes & 2 u = 0 a partial differential equation called Laplace’s equation. In Chapters 18 and 19 we will write solutions of the heat equation and Laplace’s equation under various conditions.

13.7.3 The Divergence Theorem as a Conservation of Mass Principle We will derive a model providing a physical interpretation of the divergence of a vector. Let F
x y z t be the velocity of a fluid moving in a region of 3-space  at point
x y z point in this region. Place an imaginary sphere r of radius  r about and time t. Let P0 be a 13.47. bounds a solid ball M . Let N be the unit outer normal to P0 , as in Figure r r r . We   know that r F · Nd is the flux of F out of Mr across r . If r is sufficiently small, then for a given time, & ·F
x y z t is approximated by & ·F
P0  t at all points
x y z of Mr , to within any desired tolerance. Therefore





& · F
x y z tdV ≈ & · F
P0  tdV Mr

Mr

4 = & · F
P0  t volume of Mr  = r 3 & · F
P0  t 3

N

P0 Mr

FIGURE 13.47

r

13.7 The Divergence Theorem of Gauss

571

Then

3


& · F
x y z tdV 4r 3 Mr 3

= F · Nd 4r 3   by Gauss’s theorem. Let r → 0. Then r contracts to its center P0 and this approximation becomes an equality: 3

& · F
P0  t = lim  F · Nd r→0 4r 3 r & · F
P0  t ≈

On the right is the limit, as r → 0, of the flux of F across the sphere of radius r, divided by the volume of this sphere. This is the amount per unit volume of fluid flowing out of Mr across  . Since the sphere contracts to P0 in this limit, we interpret the right side, hence also the r divergence of F at P0 , as a measure of fluid flow away from P0 . This provides a physical sense of the divergence of a vector field. In view of this interpretation, the equation




& · FdV  F · Nd = M

states that the flux of F out of M across its bounding surface exactly balances the divergence of fluid away from the points of M. This is a conservation of mass statement, in the absence of fluid produced or destroyed within M, and provides a model for the divergence theorem.

SECTION 13.7

PROBLEMS

 In each of Problems 1 through 8, evaluate  F · Nd or  div
FdV , whichever is most convenient. M  1. F = xi + yj − zk, the sphere of radius 4 about
1 1 1  2. F = 4xi − 6yj + k, the surface of the solid cylinder x2 + y2 ≤ 4, 0 ≤ z ≤ 2 (the surface includes the end caps of the cylinder)  3. F = 2yzi − 4xzj + xyk, the sphere of radius 5 about
−1 3 1  4. F = x3 i + y3 j + z3 k, the sphere of radius 1 about the origin  5. F = 4xi − zj + xk, the hemisphere x2 + y2 + z2 = 1, z ≥ 0, including the base consisting of points
x y with x2 + y2 ≤ 1  6. F =
x − yi +
y − 4xzj + xzk, the surface of the rectangular box bounded by the coordinate planes x = 0, y = 0 and z = 0 and by the planes x = 4, y = 2 and z = 3.   7. F = x2 i + y2 j + z2 k, the cone z = x2 + y2 for x2 + y2 ≤ 2, √together with the top cap consisting of points
x y 2 with x2 + y2 ≤ 2.

 8. F = x2 i − ez j + zk, the surface bounding the cylin2 2 der x + y ≤ 4, 0 ≤ z ≤ 2 (including the top and bottom caps of the cylinder)  9. Let be a smooth closed surface and F a vector field with components that are continuous with continuous  first and second partial  derivatives. through and its interior. Evaluate 
& × F · Nd. 10. Let 
x y z and 
x y z be continuous with continuous first and second partial derivatives on a  smooth closed surface and M. Suppose  its interior & = O in M. Prove that & 2 dV = 0. M 11. Show that under theconditions of Problem 10, if & = & = O, then
& 2  − & 2 dV = 0. M  12. Let be a smooth closed surface bounding an interior M. Show that 1

R · Nd volume of M = 3   where R = xi + yj + zk is a position vector for . 13. Suppose f and g satisfy Laplace’s equation in a region  M bounded by a smooth closed surface . Suppose 

f/ = g/ on . Prove that for some constant k, f
x y z = g
x y z + k for all
x y z in M.

572

13.8

CHAPTER 13

Vector Integral Calculus

The Integral Theorem of Stokes We have seen that the conclusion of Green’s theorem can be written C

F · Tds =



D


& × F · kdA

in which T is the unit tangent to C, a simple positively oriented closed curve enclosing a region D. Think of D as a flat surface with unit normal vector k and bounded by C. To generalize  to three dimensions, allow as  C to be a closed curve in 3-space, bounding a smooth surface  in Figure 13.48. Here need not be a closed surface. Let N be a unit normal to .  This raises a subtle point. At any point of there are two normal vectors, as shown in Figure 13.49. Which should we choose? In addition to this decision, we must choose a direction on C. In the plane we chose counterclockwise as positive orientation, but this has no meaning in three dimensions.   First we will give a rule for choosing a particular unit normal to at each point. If has coordinate functions x = x
u v, y = y
u v, z = z
u v, then the normal vector


z x


x y


y z i+ j+ k


u v


u v


u v  divided by its length, yields a unit normal to . The negative of this unit normal is also a unit  normal at the same point. Choose either this vector or its negative  to use as the normal to , and call it n. Whichever is chosen as n, use it at all points of . That is, do not use n at one point and –n at another. Now use the choice of n to determine an orientation or direction on C to be called the positive orientation. Referring to Figure 13.50, at any point on C, if you stand along n (that is, your head  is at the tip of n), then the positive direction on C is the one in which you walk to have over your left shoulder. The arrow shows the orientation on C obtained in this way. Although this is not a rigorous definition, it is sufficient for our purpose, without becoming enmeshed in topological details. With this direction, we say that C has been oriented coherently with n. If we had chosen the normal in the opposite direction, then we would have reached the opposite orientation on C. The choice of the normal determines the orientation on the curve. There is no intrinsic positive or negative orientation of the curve—simply an orientation coherent with the choice of normal. With this understanding, we can state Stokes’s theorem.

N

z

z

N

z N







C

C

y

y

x

x

x

FIGURE 13.48

FIGURE 13.49

FIGURE 13.50

Normals to a surface at a point.

y

13.8 The Integral Theorem of Stokes THEOREM 13.9

573

Stokes

 Let be a piecewise smooth surface bounded by a piecewise smooth curve C. Suppose a unit  normal n has been chosen on and that C is oriented coherently with n. Let F
x y z be a vector field whose  component functions are continuous with continuous first and second partial derivatives on . Then,

F · dR = 
& × F · nd C

We will write this conclusion in terms of coordinates and component functions. Let the component functions of F be, respectively, f , g and h. Then F · dR = f
x y zdx + g
x y zdy + h
x y zdz C

C

Next,

  i  & × F =  / x  f

j

/ y g

k

/ z h

          = h − g i + f − h j + g − f k 

y z

z x

x y 

The normal to the surface is given by equation (13.5) as N=


z x


y z


x y i+ j+ k


u v


u v


u v

Use this to define the unit normal n
u v =

N
u v  N
u v

Then N
u v N
u v    

f h
z x

h g
y z 1 − + − =

y z
u v

z x
u v N
u v  

g f
x y + − 

x y
u v


& × F · n =
& × F ·

Then




&

× F · nd

   

f h
z x 1

h g
y z − + −

y z
u v

z x
u v D N
u v  

g f
x y + − N
u v dudv

x y
u v

 h g 
y z  f h 
z x − + − =

y z
u v

z x
u v D 



g f
x y dudv + −

x y
u v

=





574

CHAPTER 13

Vector Integral Calculus

 in which the coordinate functions x
u v, y
u v and z
u v from are substituted into the integral, and D is set of points
u v over which these coordinate functions are defined. Keep in mind that the function to be integrated in Stokes’s theorem is
& × F · n, in which n = N/ N is a unit normal. However, in converting the surface integral to a double integral over D, using the definition of surface integral,
& × F · n must be multiplied by N
u v, with N
u v determined by equation (13.5).

EXAMPLE 13.32

  Let F
x y z = −yi + xyj − xyzk and let consist of the part of the cone z = x2 + y2 for x2 + y2 ≤ 9. We will compute both sides of the conclusion of Stokes’s theorem to illustrate the various terms and integrals involved. The cone is shown in Figure 13.51. Its boundary curve C is the circle around the top of the cone, the curve x2 + y2 = 9 in the plane z = 3. In this example the surface is described by z = S
x y. Here x and y are the parameters, varying over the disk D given by x2 + y2 ≤ 9. We can compute a normal vector

z

z i− j+k

x

y x y = − i − j + k z z

N=−

For
& × F · n in Stokes’s theorem, n is a unit normal, so compute the norm of N: # # 0 2 # x # y x y2 # N = # #− z i − z j + k# = z2 + z2 + 1 0 x2 + y2 +
x2 + y2  √ = = 2 x2 + y 2 Use the unit normal n=

1 N = √
−xi − yj + zk  N 2z

This normal is defined at all points of the cone except the origin, where there is no normal. n is an inner normal, pointing from any point on the cone into the region bounded by the cone. If we stand along n at points of C and imagine walking along C in the direction of the arrow in Figure 13.52, then the surface is over our left shoulder. Therefore this arrow orients z (0, 0, 3)

z

C: boundary curve

C: x 2  y 2  9, z  3 : z  x 2  y 2

N(3, 0, 3)

: z  x 2  y 2

y x FIGURE 13.51

y x FIGURE 13.52

13.8 The Integral Theorem of Stokes

575

C coherently with n. If we had used −n as normal vector, we would orient C in the other direction. We can parametrize C by x = 3 cos
t y = 3 sin
t z = 3

for 0 ≤ t ≤ 2

The point
3 cos
t 3 sin
t 3 traverses C in the positive direction (as determined by n) as t increases from 0 to 2. This completes the preliminary work and we can evaluate the integrals. For the line integral, F · dR = −ydx + xdy − xyz dz C

= =





C

2 0 2 0

−3 sin
t
−3 sin
t + 3 cos
t3 cos
tdt 9dt = 18

For the surface integral, first compute the curl of F:   i j k  & × F =  / x / y / z  −y x −xyz

     

= −xzi + yzj + 2k Then 1
& × F · n = √
x2 z − y2 z + 2z 2z 1 2 = √
x − y2 + 2 2 Then




& × F · nd =

= =



D







D

D


& × F · n N dxdy √ 1 √
x2 − y2 + 2 2 dx dy 2
x2 − y2 + 2dx dy

Use polar coordinates on D to write this integral as 2
3 


0

=




0



0

2
3 0

 r 2 cos2
 − r 2 sin2
 + 2 rdrd

r 3 cos
2drd +




2
3

2rdrd 0

0

2

3 1 1 4 3 = sin
2 r + 2 r 2 0 = 18 2 4 0 0 The following are two applications of Stokes’s theorem.

576

CHAPTER 13

Vector Integral Calculus

13.8.1 An Interpretation of Curl We will use Stokes’s theorem to argue for a physical interpretation of the curl operation. Think  of F
x y z as the velocity of a fluid and let P0 be any point in the fluid. Consider a disk r of radius r about P0 , with unit normal vector n and boundary circle Cr coherently oriented, as in Figure 13.53. For the disk the normal vector is constant. By Stokes’s theorem,

F · dR = 
& × F · nd Cr

r




Since R
t is/ a tangent vector to Cr , then F · R
is the tangential component of the velocity about Cr and Cr F · dR measures the circulation of the fluid about Cr . z

Boundary circle Cr

N P0

Disk  r y

x FIGURE 13.53

By choosing  r sufficiently small, & × F
x y z is approximated by & × F
P0  as closely as we like on r . Further, since n is constant,

circulation of F about Cr ≈ 
& × F
P0  · nd r

=
& × F
P0  · n
area of the disk



r

= r 2
& × F
P0  · n Therefore & × F
P0  · n ≈

1
circulation of F about Cr   r 2

As r → 0, the disk contracts to its center P0 and we obtain & × F
P0  · n = lim

r→0

1
circulation of F about Cr   r 2

Since n is normal to the plane of Cr , this equation can be read & × F
P0  · n = circulation of F per unit area in the plane normal to n Thus the curl of F is a measure of rotation of the fluid at a point. This is the reason why a fluid is called irrotational if the curl of the velocity vector is zero. For example, any conservative vector field is irrotational, because if F = &, then & × F = & ×
& = O.

13.8.2 Potential Theory in 3-Space As in the plane, a vector field F
x y z in 3-space is conservative if for some potential function , F = &. By exactly the same reasoning as applied in the two-dimensional case, we get
F · dR = 
P1  − 
P0  C

13.8 The Integral Theorem of Stokes

577

if P0 is the initial point of C and P1 the terminal point. Therefore the line integral of a conservative vector field in 3-space is independent of path. If a potential function exists, we attempt to find it by integration just as in the plane.

EXAMPLE 13.33

Let F
x y z =
yzexyz − 4xi +
xzexyz + z + cos
yj +
xyexyz + yk For this to be the gradient of a scalar field , we must have

 = yzexyz − 4x

x

 = xzexyz + z + cos
y

y

and

 = xyexyz + y

z Begin with one of these equations, say the last, and integrate with respect to z to get 
x y z = exyz + yz + k
x y where the “constant” of integration may involve the other two variables. Now we need

 = yzexyz − 4x

x

k

xyz =

e + yz + k
x y = yzexyz + 

x

x This will be satisfied if

k = −4x

x so k
x y must have the form k
x y = −2x2 + c
y Thus far 
x y z must have the appearance 
x y z = exyz + yz − 2x2 + c
y Finally, we must satisfy

 = xzexyz + z + cos
y

y 

xyz e + yz − 2x2 + c
y =

y = xzexyz + z + c

y Then c

y = cos
y and we may choose c
y = sin
y

578

CHAPTER 13

Vector Integral Calculus

A potential function is given by 
x y z = exyz + yz − 2x2 + sin
y Of course, for any number a, 
x y z + a is also a potential function for F. As in the plane, in 3-space there are vector fields that are not conservative. We would like to develop a test to determine when F has a potential function. The discussion follows that in Section 13.3.1 for the plane. As we saw in the plane, the test requires conditions not only on F, but on the set D of points on which we want to find a potential function. We define a set D of points in 3-space to be a domain if 1. about every point of D there exists a sphere containing only points of D, and 2. between any two points of D there is a path lying entirely in D. This definition is analogous to the definition made for sets in the plane. For example, the set of points bounded by two disjoint spheres is not a domain because it fails to satisfy (2), while the set of points on or inside the solid unit sphere about the origin fails to be a domain because of condition (1). The set of points
x y z with x ≥ 0, y ≥ 0 and z ≥ 0 is not a domain because it fails condition (1). For example, there is no sphere about the origin containing only points with nonnegative coordinates. The set of points
x y z with x > 0, y > 0 and z > 0 is a domain. On a domain, existence of a potential function is equivalent to independence of path. THEOREM 13.10

Let D be a domain in 3-space, and let F be continuous on D. Then path on D if and only if F is conservative.

 C

F · dR is independent of

We already know that existence of a potential function implies independence of path. For the converse, suppose C F · dR is independent of path in D. Choose any P0 in D. If P is any point of D, there is a path C in D from P0 to P, and we can define

P = F · dR C

Because this line integral depends only on the end points of any path in D, this defines a function for all
x y z in D. Now the argument used in the proof of Theorem 13.6 can be essentially duplicated to show that F = &. With one more condition on the domain D, we can derive a simple test for a vector field to be conservative. A set D of points in 3-space is simply connected if every simple closed path in D is the boundary of a piecewise smooth surface lying in D. This condition enables us to use Stokes’s theorem to derive the condition we want. THEOREM 13.11

Let D be a simply connected domain in 3-space. Let F and & × F be continuous on D. Then, F is conservative if and only if & × F = O in D. Thus, in simply connected domains, the conservative vector fields are the ones having curl zero—that is, the irrotational vector fields. In one direction, the proof is simple. If F = &, then & × F = & ×
& = O, without the requirement of simple connectivity. In the other direction, suppose & × F = O. To prove that F is conservative, it is enough to prove that C F · dR is independent of path in D. Let C and K be two paths from P0 to P1 in D.

13.8 The Integral Theorem of Stokes

579

z P1 K C

P0



y x FIGURE 13.54

Form a closed path L in D consisting of C and  −K, as in Figure 13.54. Since D is simply connected, there is a piecewise smooth surface in D having boundary C. Then

F · dR = F · dR − F · dR L

C

= so


C

K




&

F · dR =

× F · nd = 0
K

F · dR

and the line integral is independent of path, hence conservative. If G
x y = f
x yi + g
x yj, then we can think of G as a vector field in 3-space by writing G
x y = f
x yi + g
x yj + 0k Then

  i  & × G =  / x  f
x y

  j k  

g f 

/ y / z  = − k

x y g
x y 0 

so the condition & × G = O in this two-dimensional case is exactly the condition of Theorem 13.5.

SECTION 13.8

PROBLEMS

In each of 1 through 5, use Stokes’s theorem to / Problems  evaluate C F·dR or 
& ×F·Nd, whichever appears easier.  1. F = yx2 i − xy2 j + z2 k, the hemisphere x2 + y2 + 2 z = 4 z ≥ 0  2. F = xyi + yzj + xzk, the paraboloid z = x2 + y2 for 2 2 x +y ≤ 9   3. F = zi+xj+yk, the cone z = x2 + y2 for 0 ≤ z ≤ 4  the part of the paraboloid 4. F = z2 i + x2 j + y2 k, 2 2 z = 6 − x − y above the x y plane

 5. F = xyi + yzj + xyk, the part of the plane 2x + 4y + z = 8 in the first octant 6. Calculate the circulation of F =
x − yi + x2 yj + xzak counterclockwise about the circle x2 + y2 = 1. Here a is a positive constant. Hint: Use Stokes’s the orem, with any smooth surface having the circle as boundary. / 7. Use Stokes’s theorem to evaluate C F · Tds, where C is the boundary of the part of the plane x + 4y + z = 12 lying in the first octant, and F =
x − zi +
y − xj +
z − yk.

580

CHAPTER 13

Vector Integral Calculus

In each of Problems 8 through 14, let  be all of 3-space (so  is a simply connected domain). Test to see if F is conservative. If it is, find a potential function.

In each of Problems 15 through 20, evaluate the line integral of the vector field on any path from the first point to the second by finding a potential function.

8. F = 2xi − 2yj + 2zk

15. F = i − 9y2 zj − 3y2 k
1 1 1
0 3 5

9. F = i − 2j + k 10. F = yz cos
xi +
z sin
x + 1j + y sin
xk

16. F =
y cos
xz − xyz sin
xzi + x cos
xzj −x2 sin
xzk
1 0 
1 1 7

11. F =
x2 − 2i + xyzj − yz2 k

17. F = 6x2 eyz i+2x3 zeyz j+2x3 yeyz k
0 0 0
1 2 −1

12. F = exyz
1 + xyzi + x2 zj + x2 yk

18. F = −8y2 i −
16xy + 4zj − 4yk
−2 1 1
1 3 2

13. F =
cos
x + y sin
xi + x sin
xyj + k

19. F = −i + 2z2 j + 4yzk
0 0 −4
1 1 6

14. F =
2x2 + 3y2 zi + 6xyzj + 3xy2 k

20. F =
y − 4xzi + xj +
3z2 − 2x2 k
1 1 1
3 1 4

PA RT

5

CHAPTER 14 Fourier Series

CHAPTER 15 The Fourier Integral and Fourier Transforms

CHAPTER 16 Special Functions, Orthogonal Expansions and Wavelets

Fourier Analysis, Orthogonal Expansions, and Wavelets

In 1807 the French mathematician Joseph Fourier (1768–1830) submitted a paper to the Academy of Sciences in Paris. In it he presented a mathematical treatment of problems involving heat conduction. Although the paper was rejected for lack of rigor, it contained ideas that were so rich and widely applicable that they would occupy mathematicians in research to the present day. One surprising implication of Fourier’s work was that many functions can be expanded in infinite series or integrals involving sines and cosines. This revolutionary idea sparked a heated debate among leading mathematicians of the day and led to important advances in mathematics (Cantor’s work on cardinals and ordinals, orders of infinity, measure theory, real and complex analysis, differential equations), science and engineering (data compression, signal analysis), and applications undreamed of in Fourier’s day (CAT scans, PET scans, nuclear magnetic resonance). Today the term Fourier analysis refers to many extensions of Fourier’s original insights, including various kinds of Fourier series and integrals, real and complex Fourier transforms, discrete and finite transforms, and, because of their wide applicability, a variety of computer 581

programs for efficiently computing Fourier coefficients and transforms. The ideas behind Fourier series have also found important generalizations in a broad theory of eigenfunction expansions, in which functions are expanded in series of special functions (Bessel functions, orthogonal polynomials, and other functions generated by differential equations). More recently, wavelet expansions have been developed to provide additional tools in areas such as filtering and signal analysis. This part is devoted to some of these ideas and their applications.

582

CHAPTER

14

THE FOURIER SERIES OF A FUNCTION CONVERGENCE OF FOURIER SERIES FOURIER COSINE AND SINE SERIES INTEGRATION AND DIFFERENTIATION OF FOURIER SERIES THE PHASE ANGLE FORM OF A

Fourier Series

14.1

Why Fourier Series? A Fourier series is a representation of a function as a series of constants times sine and/or cosine functions of different frequencies. In order to see why such a series might be interesting, we will look at a problem of the type that led Fourier to consider them. Consider a thin bar of length , constant density and uniform cross section. Let u
x t be the temperature at time t in the cross section of the bar at x, for 0 ≤ x ≤ . In Section 13.7.2 we derived a partial differential equation for u:

u

2 u = k 2 for 0 < x <  t > 0 (14.1)

t

x in which k is a constant depending on the material of the bar. Suppose the left and right ends of the bar are kept at zero temperature, u
0 t = u
 t = 0

for t > 0

(14.2)

and that the temperature throughout the bar at time t = 0 is specified u
x 0 = f
x = x
 − x

(14.3)

Intuitively, the heat equation, together with the initial temperature distribution throughout the bar, and the information that the ends are kept at zero degrees for all time, are enough to determine the temperature distribution u
x t throughout the bar at any time. By a process that now bears his name, and which we will develop when we study partial differential equations, Fourier found functions satisfying the heat equation (14.1) and the conditions at the ends of the bar, equations (14.2), and having the form un
x t = bn sin
nxe−kn t  2

(14.4)

in which n can be any positive integer, and bn can be any real number. We will use these functions to find a function that also satisfies the condition (14.3). 583

584

CHAPTER 14

Fourier Series

Periodic phenomena have long fascinated mankind; our ancient ancestors were aware of the recurrence of phases of the moon and certain planets, the tides of lakes and oceans, and cycles in the weather. Isaac Newton’s calculus and law of gravitation enabled him to explain the periodicities of the tides, but it was left to Joseph Fourier and his successors to develop Fourier analysis, which has had profound applications in the study of natural phenomena and the analysis of signals and data.

A single choice of positive integer n0 and constant bn0 will not do. If we let u
x t = 2 bn0 sin
n0 xe−kn0 t , then we would need u
x 0 = x
 − x = bn0 sin
n0 x for 0 ≤ x ≤  an impossibility. A polynomial cannot equal a constant multiple of a sine function over 0  (or over any nontrivial interval). The next thing to try is a finite sum of the functions (14.4), say u
x t =

N 

bn sin
nxe−kn t  2

(14.5)

n=1

Such a function will still satisfy the heat equation and the conditions (14.2). To satisfy the condition (14.3), we must choose N and bn
s so that u
x 0 = x
 − x =

N 

bn sin
nx for

0 ≤ x ≤ 

n=1

But this is also impossible. A finite sum of constant multiples of sine functions cannot equal a polynomial over 0 .

14.1 Why Fourier Series?

585

At this point Fourier had a brilliant insight. Since no finite sum of functions (14.4) can be a solution, attempt an infinite series: u
x t =



bn sin
nxe−kn t  2

(14.6)

n=1

This function will satisfy the heat equation, as well as the conditions u
x 0 = u
 0 = 0. To satisfy condition (14.3) we must choose the bn
s so that u
x 0 = x
 − x =



bn sin
nx for

0 ≤ x ≤ 

(14.7)

n=1

This is quite different from attempting to represent the polynomial x
 − x by the finite trigonometric sum (14.5). Fourier claimed that equation (14.7) is valid for 0 ≤ x ≤  if the coefficients are chosen as 2
 4 1 −
−1n x
 − x sin
nxdx =  bn =  0  n3 By inserting these coefficients into the proposed solution (14.6), Fourier thus claimed that the solution of this heat conduction problem, with the given initial temperature, is u
x t =

4  1 −
−1n 2 sin
nxe−kn t   n=1 n3

The claim that  4 1 −
−1n sin
nx = x
 − x for n3 n=1 

0≤x≤

was too much for mathematicians of Fourier’s time to accept. The mathematics of this time was not adequate to proving this kind of assertion. This was the lack of rigor that led the Academy to reject publication of Fourier’s paper. But the implications were not lost on Fourier’s colleagues. There is nothing unique about x
 −x as an initial temperature distribution, and many different functions could be used. What Fourier was actually claiming  was that, for a broad class of functions f , coefficients bn could be chosen so that f
x = n=1 bn sin
nx on 0 . Eventually this and even more general claims for these series proposed by Fourier were proved. We will now begin an analysis of Fourier’s ideas and some of their ramifications.

SECTION 14.1 1. On x
 for tial 20

PROBLEMS

the same set of axes, generate a graph of 5 n 3 − x and n=1
4/
1 −
−1 /n  sin
nx 0 ≤ x ≤ . Repeat this for the par10 n 3 sums n=1
4/
1 −
−1 /n  sin
nx and n 3 n=1
4/
1 −
−1 /n  sin
nx This will give a sense of the correctness of Fourier’s intuition in asserting that x
 −x can be accurately represented by  n 3 n=1
4/
1 −
−1 /n  sin
nx on this interval.

2. Prove that a polynomial cannot be a constant multiple of sin
nx over 0 , for any positive integer n. Hint: One way is to proceed by induction on the degree of the polynomial. 3. Prove that a polynomial cannot be equal to a nonzero  sum of the form nj=0 cj sin
jx for 0 ≤ x ≤ , where
the cj s are real numbers.

CHAPTER 14

586

14.2

Fourier Series

The Fourier Series of a Function L Let f
x be defined for −L ≤ x ≤ L. For the time being, we assume only that −L f
xdx exists. We want to explore the possibility of choosing numbers a0  a1      b1  b2     such that  nx  nx  1 f
x = a0 + an cos + bn sin 2 L L n=1

(14.8)

for −L ≤ x ≤ L. We will see that this is sometimes asking too much, but that under certain conditions on f it can be done. However, to get started, we will assume the best of all worlds and suppose for the moment that equation (14.8) holds. What does this tell us about how to choose the coefficients? There is clever device used to answer this question, which was known to Fourier and others of his time. We will need the following elementary lemma.

LEMMA 14.1

Let n and m be nonnegative integers. Then 1.


L −L

cos

 nx L

sin

 mx L

dx = 0

2. If n = m, then


L −L

cos

 nx L

cos

 mx L

dx =




L −L

sin

 nx L

sin

 mx L

dx = 0

3. If n = 0, then


L −L

cos2

 nx L

dx =




L −L

sin2

 nx L

dx = L

The lemma is proved by straightforward integration. Now, to find a0 , integrate the series (14.8) term by term (supposing for now that we can do this):


1
L f
xdx = a0 dx 2 −L −L
L
L  nx  nx  + an dx + bn dx cos sin L L −L −L n=1 L

All of the integrals on the right are zero and this equation reduces to


L

−L

f
xdx = La0 

Therefore a0 =

1
L f
xdx L −L

14.2 The Fourier Series of a Function

587

Next, we will determine ak for any positive integer k. Multiply equation (14.8) by cos
kx/L and integrate each term of the resulting series to get    
L 1
L kx kx dx = a0 dx f
x cos cos L 2 L −L −L    
L
L  nx  nx  kx kx + an cos cos cos sin dx + bn dx L L L L −L −L n=1 L By the lemma, all of the integrals on the right are zero except for −L cos
kx/L cos
kx/Ldx, which occurs when n = k, and in this case this integral equals L. The right side of this equation therefore collapses to just one term, and the equation becomes  
L kx f
x cos dx = ak L L −L whereupon   1
L kx dx ak = f
x cos L −L L To determine bk , return to equation (14.8). This time multiply the equation by sin
kx/L and integrate each term to get    
L kx kx 1
L f
x sin sin dx = a0 dx L 2 L −L −L
L
L  nx  kx   nx  kx   sin sin + an cos sin dx + bn dx L L L L −L −L n=1 Again, by the lemma, all terms on the right are zero except for when n = k, and this equation reduces to  
L kx f
x sin dx = bk L L −L Therefore bk =

L −L

sin
nx/L sin
kx/Ldx

  kx 1
L f
x sin dx L −L L

We have now “solved” for the coefficients in the trigonometric series expansion 14.8. Of course, this analysis is flawed by the interchange of series and integrals, which is not always justified. However, the argument does tell us how the constants should be chosen, at least under certain conditions, and suggests the following definition. DEFINITION 14.1

Fourier Coefficients and Series

Let f be a Riemann integrable function on −L L. 1. The numbers an =

 nx 1
L dx, f
x cos L −L L

for

n = 0 1 2   

588

CHAPTER 14

Fourier Series

and bn =

 nx 1
L dx f
x sin L −L L

for n = 1 2 3   

are the Fourier coefficients of f on −L L. 2. The series  nx  nx  1 a0 + an cos + bn sin 2 L L n=1

is the Fourier series of f on −L L when the constants are chosen to be the Fourier coefficients of f on −L L.

EXAMPLE 14.1

Let f
x = x for − ≤ x ≤ . We will write the Fourier series of f on − . The coefficients are: 1
 xdx = 0 a0 =  − 1
 x cos
nxdx an =  −

 1 x sin
nx = 0 = 2 cos
nx + n n − and

1
 x sin
nxdx  −

 1 x cos
nx = 2 sin
nx − n n −

bn =

2 2 = − cos
n =
−1n+1  n n since cos
n =
−1n if n is an integer. The Fourier series of x on −  is  1 2 2 2
−1n+1 sin
nx = 2 sin
x − sin
2x + sin
3x − sin
4x + sin
5x − · · ·  3 2 5 n=1 n

In this example the constant term and cosine coefficients are all zero, and the Fourier series contains only sine terms.

EXAMPLE 14.2

Let

 f
x =

0 x

for − 3 ≤ x < 0  for 0 ≤ x ≤ 3

14.2 The Fourier Series of a Function

589

Here L = 3 and the Fourier coefficients are: 1
3 1
3 3 f
xdx = xdx =  3 −3 3 0 2
 3 1 nx dx an = f
x cos 3 −3 3  nx 1
3 = x cos dx 3 0 3  nx 3  nx 3 x = 2 2 cos + sin n 3 n 3 0 a0 =

=

3 n2  2


−1n − 1

and  nx  nx 1
3 1
3 dx = dx f
x sin x sin 3 −3 3 3 0 3  nx 3  nx 3 x = 2 2 sin − cos n 3 n 3 0

bn =

=

3
−1n+1  n

The Fourier series of f on −3 3 is   nx  nx  3  3 3 n n+1 + +
−1


−1 − 1 cos sin  4 n=1 n2  2 3 n 3

L L Even when f
x is fairly simple, −L f
x cos
nx/Ldx and −L f
x sin
nx/Ldx can involve considerable labor if done by hand. Use of a software package to evaluate definite integrals is highly recommended. In these examples, we wrote the Fourier series of f , but did not claim that it equalled f
x. For most x it is not obvious what the sum of the Fourier series is. However, in some cases it is obvious that the series does not equal f
x. Consider again f
x = x on −  in Example 14.1. At x =  and at x = −, every term of the Fourier series is zero, even though f
 =  and f
− = −. Even for very simple functions, then, there may be points where the Fourier series does not converge to f
x. Shortly we will determine the sum of the Fourier series of a function. Until this is done, we do not know the relationship between the Fourier series and the function itself.

14.2.1

Even and Odd Functions

Sometimes we can save some work in computing Fourier coefficients by observing special properties of f
x.

590

CHAPTER 14

Fourier Series

DEFINITION 14.2 Even Function

f is an even function on −L L if f
−x = f
x for −L ≤ x ≤ L. Odd Function

f is an odd function on −L L if f
−x = −f
x for −L ≤ x ≤ L.

For example, x2 , x4 , cos
nx/L and e−x are even functions on any interval −L L. Graphs of y = x2 and y = cos
5x/3 are given in Figure 14.1. The graph of such a function for −L ≤ x ≤ 0 is the reflection across the y-axis of the graph for 0 ≤ x ≤ L (Figure 14.2). The functions x, x3 , x5 and sin
nx/L are odd functions on any interval −L L. Graphs of y = x, y = x3 and y = sin
5x/2 are shown in Figure 14.3. The graph of an odd function for −L ≤ x ≤ 0 is the reflection across the vertical axis, and then across the horizontal axis, of the graph for 0 ≤ x ≤ L (Figure 14.4). If f is odd, then f
0 = 0, since f
−0 = f
0 = −f
0 Of course, “most” functions are neither even nor odd. For example, f
x = ex is not even or odd on any interval −L L.

y

y


3
2

4


2

8 6 4 2


3


1

1

2

3

x

FIGURE 14.1 Graphs of even functions y = x2 and y = cos
5x/3.


1 0
4
8

1

2

x

3


12 FIGURE 14.2 Graph of a typical even function, symmetric about the y axis.

y y 8 8

4

4
2


1

1
6

FIGURE 14.3 Graphis of odd functions y = x, y = x3 , and y = sin
5x/2.

2

x


2


1


4

1

2

x


8 FIGURE 14.4 Graph of a typical odd function, symmetric through the origin.

14.2 The Fourier Series of a Function

591

Even and odd functions behave as follows under multiplication: even · even = even, odd · odd = even, and even · odd = odd. For example, x2 cos
nx/L is an even function (product of two even functions); x2 sin
nx/L is odd (product of an even function with an odd function); and x3 sin
nx/L is even (product of two odd functions). Now recall from calculus that
L f
xdx = 0 if f is odd on −L L −L

and


L −L

f
xdx = 2




L

f
xdx

if

f is even on −L L.

0

These integrals are suggested by Figures 14.2 and 14.4. In Figure 14.4, f is odd on −L L, and the “area” bounded by the graph and the horizontal axis for −L ≤ x ≤ 0 is exactly  L the negative of that bounded by the graph and the horizontal axis for 0 ≤ x ≤ L. This makes −L f
xdx = 0. In Figure 14.2, where f is even, the area to the leftof the vertical axis, for −L ≤ x ≤ 0, equals L L that to the right, for 0 ≤ x ≤ L, so −L f
xdx = 2 0 f
xdx. One ramification of these ideas for Fourier coefficients is that, if f is an even or odd function, then some of the Fourier coefficients can be seen immediately to be zero, and we need not carry out the integrations explicitly. We saw this in Example 14.1 with f
x = x, which is an odd function on − . There we found that the cosine coefficients were all zero, since x cos
nx is an odd function.

EXAMPLE 14.3

We will find the Fourier series of f
x = x4 on −1 1. Since f is an even function, x4 sin
nx is odd and we know immediately that all the sine coefficients bn are zero. For the other coefficients, compute: a0 =




1 −1

x4 dx = 2




1 0

2 x4 dx =  5

and an =




1 −1

=2




x4 cos
nxdx 1

0

x4 cos
nxdx = 8

n2  2 − 6
−1n   4 n4

The Fourier series of x4 on −1 1 is 1  n2  2 − 6 + 8
−1n cos
nx 5 n=1  4 n4

592

CHAPTER 14

Fourier Series

To again make the point about convergence, notice that f
0 = 0 in this example, but the Fourier series at x = 0 is 1  n2  2 − 6 + 8
−1n  5 n=1  4 n4 It is not clear whether or not this series sums to the function value 0.

EXAMPLE 14.4

Let f
x = x3 for −4 ≤ x ≤ 4. Because f is odd on −4 4, its Fourier cosine coefficients are all zero. Its Fourier sine coefficients are  nx 1
4 3 bn = dx x sin 4 −4 4  nx 1
4 3 n2  2 − 6 = dx =
−1n+1 128 3 3  x sin 2 0 4 n The Fourier series of x3 on −4 4 is 


−1n+1 128

n=1

 nx n2  2 − 6  sin n3  3 4

We will make use of this discussion later, so here is a summary of its conclusions: If f is even on −L L, then its Fourier series on this interval is  nx  1 a0 + an cos  (14.9) 2 L n=1 in which

 nx 2
L dx for n = 0 1 2     f
x cos L 0 L If f is odd on −L L, then its Fourier series on this interval is  nx   bn sin L n=1 an =

where bn =

SECTION 14.2

 nx 2
L dx f
x sin L 0 L

PROBLEMS

In each of Problems 1 through 12, find the Fourier series of the function on the interval. 1. f
x = 4 −3 ≤ x ≤ 3 2. f
x = −x −1 ≤ x ≤ 1

for n = 1 2    

3. f
x = cosh
x = 1 ≤ x ≤ 1 4. f
x = 1 − x  −2 ≤ x ≤ 2  −4 for −  ≤ x ≤ 0 5. f
x = 4 for 0 < x ≤ 

(14.10)

(14.11)

(14.12)

14.3 Convergence of Fourier Series 6. f
x = sin
2x − ≤ x ≤ 

12. f
x =

7. f
x = x2 − x + 3 −2 ≤ x ≤ 2  −x for − 5 ≤ x < 0 8. f
x = 1 + x2 for 0 ≤ x ≤ 5  1 for −  ≤ x < 0 9. f
x = 2 for 0 ≤ x ≤  10. f
x = cos
x/2 − sin
x − ≤ x ≤  11. f
x = cos
x −3 ≤ x ≤ 3

14.3

 1−x 0

593

for − 1 ≤ x ≤ 0 for 0 < x ≤ 1

13. Suppose f and g are integrable on −L L and that f
x = g
x except for x = x0 , a given point in the interval. How are the Fourier series of f and g related? What does this suggest about the relationship between a function and its Fourier series on an interval? L 14. Prove that −L f
xdx = 0 if f is odd on −L L. L L 15. Prove that −L f
xdx = 2 0 f
xdx if f is even on

−L L.

Convergence of Fourier Series It is one thing to be able to write the Fourier coefficients of a function f on an interval L L

−L L. This requires only existence of −L f
x cos
nx/Ldx and −L f
x sin
nx/Ldx. It is another issue entirely to determine whether the resulting Fourier series converges to f
x—or even that it converges at all! The subtleties of this question were dramatized in 1873 when the French mathematician Paul Du Bois-Reymond gave an example of a function which is continuous on
− , but whose Fourier series fails to converge at any point of this interval. However, the obvious utility of Fourier series in solving partial differential equations led in the nineteenth century to an intensive effort to determine their convergence properties. About 1829, Peter Gustav Lejeune-Dirichlet gave conditions on the function f which were sufficient for convergence of the Fourier series of f . Further, Dirichlet’s theorem actually gave the sum of the Fourier series at each point, whether or not this sum is f
x. This section is devoted to conditions on a function that enable us to determine the sum of its Fourier series on an interval. These conditions center about the concept of piecewise continuity.

DEFINITION 14.3

Piecewise Continuous Function

Let f
x be defined on a b, except possibly at finitely many points. Then f is piecewise continuous on a b if: 1. f is continuous on a b except perhaps at finitely many points. 2. Both limx→a+ f
x and limx→b− f
x exist and are finite. 3. If x0 is in
a b and f is not continuous at x0 , then limx→x0 + f
x and limx→x0 − f
x exist and are finite.

Figures 14.5 and 14.6 shows graphs of typical piecewise continuous functions. At a point of discontinuity (which we assume are finite in number), the function must have finite one-sided limits. This means that the graph experiences at worst a finite gap at a discontinuity. Points where these occur are called jump discontinuities of the function. As an example of a simple function that is not piecewise continuous, let  0 for x = 0 f
x =  1/x for 0 < x ≤ 1

594

CHAPTER 14

Fourier Series y y x x FIGURE 14.5 A piecewise continuous function.

FIGURE 14.6

Graph of a typical piecewise continuous function.

Then f is continuous on
0 1, and discontinuous at 0. However, limx→0+ f
x = , so the discontinuity is not a finite jump discontinuity, and f is not piecewise continuous on 0 1.

EXAMPLE 14.5

Let

f
x =

⎧ 5 ⎪ ⎪ ⎪ ⎪ ⎨x

for x = − for −  < x < 1

⎪ 1 − x2 ⎪ ⎪ ⎪ ⎩ 4

for 1 ≤ x < 2



for 2 ≤ x ≤ 

A graph of f is shown in Figure 14.7. This function is discontinuous at −, and lim f
x = −

x→−+

f is also discontinuous at 1, interior to − , and lim f
x = 1

and

x→1−

lim f
x = 0

x→1+

Finally, f is discontinuous at 2, and lim f
x = −3

and

x→2−

lim f
x = 4

x→2+

At each point of discontinuity interior to the interval, the function has finite one sided limits from both sides. At the point of discontinuity at the end point −, the function has a finite limit y 5 4 1


1 2

x

FIGURE 14.7 Graph of the function of Example 14.5.

14.3 Convergence of Fourier Series

595

from within the interval. In this example, the other end point is not an issue, as f is continuous (from the left) there. Therefore f is piecewise continuous on − . We will use the following notation for left and right limits of a function at a point: f
x0 + = lim f
x and f
x0 − = lim f
x x→x0 +

x→x0 −

In Example 14.5, f
1− = 1 and f
1+ = 0 and f
2− = −3 and f
2+ = 4 At the end points of an interval, we can still use this notation, except at the left end point we consider only the right limit (from inside the interval), and at the right end point we use only the left limit (again, so that the limit is taken from within the interval). Again referring to Example 14.5, f
−+ = −

DEFINITION 14.4

and

f
− = 4

Piecewise Smooth Function

f is piecewise smooth on a b if f and f
are piecewise continuous on a b.

A piecewise smooth function is therefore one that is continuous except possibly for finitely many jump discontinuities, and has a continuous derivative at all but finitely many points, where the derivative may not exist but must have finite one-sided limits.

EXAMPLE 14.6

Let

⎧ 1 ⎪ ⎨ f
x = −2x ⎪ ⎩ 9e−x

for − 4 ≤ x < 1 for 1 ≤ x < 2



for 2 ≤ x ≤ 3

Figure 14.8 shows a graph of f . The function is continuous except for finite jump discontinuities at 1 and 2. Therefore f is piecewise continuous on −4 3. The derivative of f is ⎧ 0 for − 4 < x < 1 ⎪ ⎨
for 1 < x < 2 f
x = −2  ⎪ ⎩ −x −9e for 2 < x < 3 The derivative is continuous on −4 3 except at the points of discontinuity 1 and 2 of f , where f

x does not exist. However, at these points f

x has finite one-sided limits. Thus f
is piecewise continuous on −4 3, so f is piecewise smooth. As suggested by Figure 14.8, a piecewise smooth function is one that has a continuous tangent at all but finitely many points. We will now state our first convergence theorem.

596

CHAPTER 14

Fourier Series y 1
4
3
2
1
1

1

2

3

4

x


2
3
4

Graph of ⎧ 1 for −4 ≤ x < 1 ⎪ ⎨ for 1 ≤ x < 2 f
x = −2x ⎪ ⎩ −x 9e for 2 ≤ x ≤ 3 FIGURE 14.8

THEOREM 14.1

Convergence of Fourier Series

Let f be piecewise smooth on −L L. Then for −L < x < L, the Fourier series of f on

−L L converges to 1
f
x+ + f
x−  2 This means that, at each point between −L and L, the function converges to the average of its left and right limits. If f is continuous at x, then these left and right limits both equal f
x, so the Fourier series converges to the function value at x. If f has a jump discontinuity at x, then the Fourier series may not converge to f
x, but will converge to the point midway between the ends of the gap in the graph at x (Figure 14.9). y f(x) 1 2

( f (x)  f (x
))
L

x x f(x
)

L

FIGURE 14.9 Convergence of a Fourier series at a jump discontinuity.

EXAMPLE 14.7

Let

⎧ ⎪ 5 sin
x ⎪ ⎪ ⎪ ⎪ ⎪ 4 ⎨ f
x = x2 ⎪ ⎪ ⎪ 8 cos
x ⎪ ⎪ ⎪ ⎩4x

for for for for for

− 2 ≤ x < −/2 x = −/2 − /2 < x < 2 2≤x 0 for A < t < B, 2. 
A = a and 
B = b, and 3. #
p = $

p for A ≤ t ≤ B. If we think of t = 
p, then $
t = #
p. The curves have the same initial and terminal points and the same graph and orientation, but $
t moves along the graph as t varies from a to b, while #
p moves along the same graph in the same direction as p varies from A to B. Informally, two curves are equivalent if one is just a reparametrization of the other.
1  3i 2
1

y 1 1

x

The join of $1
t = eit for 0 ≤ t ≤ , with $2
t = −1 + it for 0 ≤ t ≤ 3.

FIGURE 22.7

EXAMPLE 22.5

Let $
t = t2 − 2ti for 0 ≤ t ≤ 1 and #
p = sin2
p − 2 sin
pi for 0 ≤ p ≤ /2 Both of these curves have the same graph (Figure 22.8), extending from initial point 0 to terminal point 1 − 2i. Let t = 
p = sin
p for 0 ≤ p ≤ /2 Then  is a differentiable, increasing function that takes 0 /2 onto 0 1. Further, for 0 ≤ p ≤ /2, $
sin
p = sin2
p − 2 sin
pi = #
p These curves are therefore equivalent. Informally, we will often describe a curve geometrically and speak of a curve and its graph interchangeably. When this is done, it is important to keep track of the orientation along the curve, and whether or not the curve is closed. For example, suppose $ is the straight line from 1 + i to 3 + 3i (Figure 22.9). This gives the graph and its orientation, and from this we can find $. Since the graph is the segment of the straight line from
1 1 to
3 3, the coordinate functions are x = t y = t for 1 ≤ t ≤ 3 Then $
t = x
t + y
ti =
1 + it for 1 ≤ t ≤ 3

22.1 Curves in the Plane 0

0.2

0.4

0.6

0.8

1.0

979

x


0.5

y 3  3i


1.0

y


1.5 3i


2.0

1i

x

i

(t)  i  2e it, 0  t  2 2i x

y FIGURE 22.8

for 0 ≤ t ≤ 1.

Graph of $
t = t2 − 2it

FIGURE 22.9

FIGURE 22.10

Directed line from 1 + i to 3 + 3i.

is one representation of the curve that has been described. There are of course other, equivalent representations. As another example, suppose # is the quarter-circle of radius 2 about i, from 2 + i to 3i (Figure 22.10). Again, we have given the graph and its orientation. Using polar coordinates centered at i =
0 1, we can write the coordinate functions x
t = 2 cos
t y
t = 1 + 2 sin
t for 0 ≤ t ≤ /2. As a function, this curve can be written #
t = 2 cos
t + 2i sin
t + i = i + 2eit for 0 ≤ t ≤ /2 Other, equivalent, representations can also be used. Finally, we will often make statements such as “f is continuous on $,” by which we mean that f is a complex function that is continuous at all points on the graph of $. And when we refer to “z on $,” we mean a complex number z lying on the graph of $. Curves are the objects over which we integrate complex functions. We will now define this integral.

SECTION 22.1

PROBLEMS

In each of Problems 1 through 10, graph the curve, determine its initial and terminal points, whether it is closed or not closed, whether or not it is simple, and the tangent to the curve at each point where the tangent exists. This tangent may be expressed as a vector or as a complex function.

4. $
t = 3 cos
t + 5 sin
ti for 0 ≤ t ≤ 2 5. )
t = 3 cos
t + 5 sin
ti for 0 ≤ t ≤ 4 6. ,
t = 4 sin
t − 2 cos
ti for − ≤ t ≤ /2 7. +
t = t − t2 i for −2 ≤ t ≤ 4

1. $
t = 4 − 2i + 2eit for 0 ≤ t ≤ 

1 8. #
t =
2t + 1 − t2 i for −3 ≤ t ≤ −1 2

2. $
t = ie2it for 0 ≤ t ≤ 2

9. $
t = cos
t − 2 sin
2ti for 0 ≤ t ≤ 2

3. $
t = t + t i for 1 ≤ t ≤ 3 2

10.
t = t2 − t4 i for −1 ≤ t ≤ 1

CHAPTER 22

980

22.2

Complex Integration

The Integral of a Complex Function We will define the integral of a complex function in two stages, beginning with the special case that f is a complex function defined on an interval a b of real numbers. An example of such a function is f
x = x2 + sin
xi for 0 ≤ x ≤ . It is natural to integrate such a function as


 1 f
xdx = x2 dx + i sin
xdx =  3 + 2i 3 0 0 0 This is the model we follow in general for such functions.

DEFINITION 22.1

Let f  a b →  be a complex function. Let f
x = u
x + iv
x for a ≤ x ≤ b. Then
b
b
b f
xdx = u
xdx + i v
xdx a

a

a

Both of the integrals on the right are Riemann integrals of real-valued functions over a b.

EXAMPLE 22.6

Let f
x = x − ix2 for 1 ≤ x ≤ 2. Then
2
2
2 3 7 f
xdx = xdx − i x2 dx = − i 2 3 1 1 1

EXAMPLE 22.7

Let f
x = cos
2x + i sin
2x for 0 ≤ x ≤ /4. Then
/4
/4
/4 1 1 f
xdx = cos
2xdx + i sin
2xdx = + i 2 2 0 0 0 In the last example, it is tempting to let f
x = e2ix and adapt the fundamental theorem of calculus to complex functions to obtain


/4
/4  1 2ix /4 1  i/2 e e −1 f
xdx = e2ix dx = = 2i 2i 0 0 0 =

1   1 1
cos
 + i sin
 − 1 =
−1 + i =
1 + i 2i 2 2 2i 2

We will justify this calculation shortly. We can now define the integral of a complex function over a curve in the plane.

22.2 The Integral of a Complex Function

981

DEFINITION 22.2

Let f be a complex function. Let $  a b →  be a smooth curve in the plane. Assume that f is continuous at all points on $. Then the integral of f over $ is defined to be
b
f
zdz = f
$
t$

tdt $

a

Since z = $
t on the curve, this integral is often written

b f
zdz = f
z
tz

tdt $ a  This formulation has the advantage of suggesting the way $ f
zdz is evaluated – replace z with z
t on the curve, let dz = z

tdt, and integrate over the interval a ≤ t ≤ b.

EXAMPLE 22.8

 Evaluate $ zdz if $
t = eit for 0 ≤ t ≤ . The graph of $ is the upper half of the unit circle, oriented counterclockwise from 1 to −1 (Figure 22.11). On $, z
t = eit and z

t = ieit . Further, f
z
t = z
t = e−it because t is real. Then


 f
zdz = e−it ieit dt = i dt = i 0

$

0

EXAMPLE 22.9

 Evaluate # z2 dz if #
t = t + it for 0 ≤ t ≤ 1. The graph of # is the straight line segment from the origin to
1 1, as shown in Figure 22.12. On the curve, z
t =
1 + it. Since f
z = z2 , f
z
t =
z
t2 =
1 + i2 t2 = 2it2 and z

t = 1 + i Then


#

z2 dz =




1 0

2it2
1 + idt =
−2 + 2i




1 0

2 t2 dt =
−1 + i 3

y z(t)  eit, 0  t  
1

1

FIGURE 22.11

x

y

1i z(t)  t  it, 0  t  1 x

FIGURE 22.12

982

CHAPTER 22

Complex Integration

EXAMPLE 22.10

 Evaluate $ z Re
zdz if $
t = t − it2 for 0 ≤ t ≤ 2. Here f
z = z Re
z, and on this curve, z
t = t − it2 , so f
z
t = z
t Re
z
t =
t − it2 t = t2 − it3  Further, z

t = 1 − 2it so
$

z Re
zdz = =




2 0




0

2


t2 − it3 
1 − 2itdt =
t2 − 2t4 dt − 3i


0

2




2 0


t2 − 3it3 − 2t4 dt

t3 dt = −

152 − 12i 15

We will show that the integrals of a function over equivalent curves are equal. This is important because we can parametrize a given curve infinitely many different ways, and this should not change the value of the integral of a given function over the curve.

THEOREM 22.1

Let $ and # be equivalent curves and let f be continuous on their graph. Then
$

f
zdz =


f
zdz #

Proof Suppose $  a b →  and #  A B → . Because these curves are equivalent, there is a continuous function  with positive derivative on A B such that 
A = a and 
B = b, and #
p = $

p for A ≤ p ≤ B. By the chain rule,

#

p = $


p

p Then
#

f
zdz =




B A

f
#
p#

pdp =




B

f
$

p$


p

pdp

A

Let s = 
p. Then s varies from a to b as p varies from A to B. Continuing from the last equation, we have
#

f
zdz =




b a

f
$
s$

sds =


f
zdz $

Thus far we can integrate only over smooth curves. We can extend the definition to an integral over piecewise smooth curves by adding the integrals over the components of the join.

22.2 The Integral of a Complex Function

983

DEFINITION 22.3

Let $ = $1 ⊕ $2 ⊕ · · · ⊕ $n be a join of smooth curves. Let f be continuous on each $j . Then
n
 f
zdz = f
zdz $

j=1 $j

EXAMPLE 22.11

Let $1
t = 3eit for 0 ≤ t ≤ /2, and let $2
t = t2 + 3i
t + 1 for 0 ≤ t ≤ 1. $1 is the quarter circle of radius 3 about the origin, extending counterclockwise from 3 to 3i, and $2 is the part of the parabola x =
y − 32 /9 from 3i to 1 + 6i. Figure 22.13 shows a graph of $ = $1 ⊕ $2 . We will evaluate $ Im
zdz.

y 1  6i 2 3i

On $1  write z
t = 3eit = 3 cos
t + 3i sin
tThen
/2
/2
Im
zdz = Im
z
tz

tdt = 3 sin
t −3 sin
t + 3i cos
tdt

1 3

x

0

$1

0

= −9

FIGURE 22.13


0

/2

sin2
tdt + 9i




/2 0

9 9 sin
t cos
tdt = −  + i 4 2

On $2  z
t = t2 + 3i
t + 1 and z

t = 2t + 3i so

1 Im
zdz = Im t2 + 3i
t + 1 2t + 3idt 0

$2

= =





1 0 1 0

3
t + 1
2t + 3idt = 6t2 + 6tdt + 9i


0

1


1 0

 6t2 + 6t + 9it + 9i dt


t + 1dt = 5 +

27 i 2

Then
$

22.2.1

9 27 9 9 f
zdz = −  + i + 5 + i = 5 −  + 18i 4 2 2 2

The Complex Integral in Terms of Real Integrals

It is possible to think of the integral of a complex function over a curve as a sum of line integrals of real-valued functions of two real variables over the curve. Let f
z = u
x y + iv
x y and, on the curve $, suppose z
t = x
t + iy
t for a ≤ t ≤ b. Now f
z
t = u
x
t y
t + iv
x
t y
t and z

t = x

t + iy

t

984

CHAPTER 22

Complex Integration

so f
z
tz

t = u
x
t y
t + iv
x
t y
t x

t + iy

t = u
x
t y
tx

t − v
x
t y
ty

t + i v
x
t y
tx

t + u
x
t y
ty

t Then
$

f
zdz =




b a

+i

u
x
t y
tx

t − v
x
t y
ty

t dt
a

b

v
x
t y
tx

t + u
x
t y
ty

tdt

In the notation of real line integrals,


f
zdz = udx − vdy + i vdx + udy $

$

(22.1)

$

This formulation allows a perspective that is sometimes useful in developing properties of complex integrals.

EXAMPLE 22.12

 Evaluate $ iz2 dz if $
t = 4 cos
t + i sin
t for 0 ≤ t ≤ /2. Figure 22.14 shows a graph of $, which is part of the ellipse

To evaluate

 $

x2 + y2 = 1 16 iz2 dz in terms of real line integrals, first compute f
z = iz2 = −2xy + i
x2 − y2  = u + iv

where u
x y = −2xy and v
x y = x2 − y2 

y 1.0 0.8 0.6 0.4 0.2 0

1

FIGURE 22.14

for 0 ≤ t ≤ /2.

2

3

4

x = 4 cos
t, y = sin
t

x

22.2 The Integral of a Complex Function

985

On the curve, x
t = 4 cos
t and y
t = sin
t. Now equation (22.1) gives us
/2 

/2 
−8 cos
t sin
t
−4 sin
tdt − 16 cos2
t − sin2
t cos
tdt iz2 dz = 0

$

+i =




0

/2 



16 cos2
t − sin2
t
−4 sin
tdt +

0




/2


−8 cos
t sin
t cos
tdt 0

1 64 − i 3 3

We will have an easier way of evaluating simple line integrals such as have more properties of complex integrals.

22.2.2

$

iz2 dz, when we

Properties of Complex Integrals

We will develop some properties of

THEOREM 22.2





$

f
zdz.

Linearity

Let $ be a piecewise smooth curve and let f and g be continuous on $. Let  and  be complex numbers. Then



f
z + g
zdz =  f
zdz +  g
zdz $

$

$

This conclusion is certainly something we expect of anything called an integral. The result extends to arbitrary finite sums:
 n $ j=1

j fj
zdz =

n  j=1


j

$

fj
zdz

Orientation plays a significant role in the complex integral, because it is an intrinsic part of the curve over which the integral is taken. Suppose $  a b →  is a smooth curve, as typically shown in Figure 22.15. The arrow indicates orientation. We can reverse this orientation by defining the new curve $r
t = $
a + b − t for a ≤ t ≤ b $r is a smooth curve having the same graph as $. However, $r
a = $
b and $r
b = $
a $r starts where $ ends, and $r ends where $ begins. The orientation has been reversed. We claim that reversing orientation changes the sign of the integral. y  (b) FIGURE 22.15

y

(a) x

r r (a)  (b)

Reversing orientation on a curve.

r (b)  (a) x

986

CHAPTER 22

Complex Integration

THEOREM 22.3

Reversal of Orientation

Let $  a b → C be a smooth curve. Let f be continuous on $. Then

f
zdz = − f
zdz $

Proof

$r

Let u = a + b − t. By the chain rule, d d $r
t = $
a + b − t = $

uu

t = −$

u = −$

a + b − t dt dt

Then


$r

f
zdz =




b a

f
$r
t$r

tdt = −




b a

f
$
a + b − t$

a + b − tdt

Now change variables by putting s = a + b − t. Then
a
b

f
zdz = − f
$
s$

s
−1ds = − f
$
s$

sds = − f
zdz $r

b

a

$

We need not actually define $r to reverse orientation in specific integrals — just integrate from b to a instead of from a to b. This reverses the roles of the initial and terminal points and hence the orientation. Or we can integrate from a to b and take the negative of the result. We will next state a complex version of the fundamental  theorem of calculus. It states that, if f has a continuous antiderivative F , then the value of $ f
zdz is the value of F at the terminal point of $, minus the value of F at the initial point. THEOREM 22.4

Let f be continuous on an open set G, and suppose F

z = f
z for z in G. Let $  a b → G be a smooth curve in G. Then
b f
zdz = F
$
b − F
$
a a

Proof

With $
t = z
t = x
t + iy
t, and F
z = U
x y + iV
x y,
b
b
b d
F
z
tdt f
zdz = f
z
tz

tdt = F

z
tz

tdt = a a a dt $
b d
b d U
x
t y
tdt + i V
x
t y
tdt = a dt a dt

Now we can apply the fundamental theorem of calculus to the two real integrals on the right to obtain
f
zdz = U
x
b y
b + iV
x
b y
b − U
x
a y
a + iV
x
a y
a $

= F
x
b y
b − F
x
a y
a = F
$
b − F
$
a

EXAMPLE 22.13

We will compute

 $


z2 + izdz if $
t = t5 − t cos
ti for 0 ≤ t ≤ 1.

22.2 The Integral of a Complex Function

987

1 This is an elementary but tedious calculation if done by computing 0 f
z
tz

tdt. However, if we let G be the entire complex plane, then G is open, and F
z = z3 /3 + iz2 /2 satisfies F

z = f
z. The initial point of $ is $
0 = 0 and the terminal point is $
1 = 1 − cos
1i. Therefore

z2 + izdz = F
$
1 − F
$
0 = F
1 − cos
1i − F
0 $

  1 1 1 i
1 − cos
1i + i  =
1 − cos
1i3 +
1 − cos
1i2 =
1 − cos
1i2 3 2 3 2

 One ramification of Theorem 22.4 is that, under the given conditions, the value of $ f
zdz depends only on the initial and terminal points of the curve. If # is also a smooth curve in G having the same initial point as $ and the same terminal point as $, then

f
zdz = f
zdz $

#

This is called independence of path, about which we will say more later. Another consequence is that, if $ is a closed curve in G, then the initial and terminal points coincide and
f
zdz = 0 $

We will consider this circumstance in more detail when we consider Cauchy’s theorem. The next result is used in bounding the magnitude of an integral, as we sometimes need to do in making estimates in equations or inequalities. THEOREM 22.5

Let $  a b → C be a smooth curve and let f be continuous on $. Then 

b   f
z
t z

t dt  f
zdz ≤ $

a

If, in addition, there is a positive number M such that f
z ≤ M for all z on $, then 
    f
zdz ≤ ML $

where L is the length of $. Proof



Write the complex number

$

f
zdz in polar form:
f
zdz = rei  $

Then r = e−i


$

f
zdz = e−i




b

f
z
tz

tdt

a

Since r is real,




b
b

 r = Re
r = Re e−i f
z
tz

tdt = Re e−i f
z
tz

t dt a

a

Now for any complex number w, Re
w ≤ w. Therefore 

  Re e−i f
z
tz

t ≤ e−i f
z
tz

t = f
z
tz

t 

988

CHAPTER 22

Complex Integration

  since e−i  = 1 for  real. Then 
b 
b 

b    
 f
z
tz

t dt = f
z
t z

t dt  f
zdz = r =  f
z
tz
tdt ≤ a

$

a

a

as we wanted to show. If now f
z ≤ M on $, then 

b
b   f
zdz f
z
t z

t dt ≤ M z

t dt  ≤ $

a

a

If $
t = x
t + iy
t, then z

t = x

t + iy

t = so




x

t2 +
y

t2 



b  
x

t2 +
y

t2 dt = ML  f
zdz ≤ M a

$

EXAMPLE 22.14

  We will obtain a bound on  $ eRe
z dz, where $ is the circle of radius 2 about the origin, traversed once counterclockwise. On $ we can write z
t = 2 cos
t + 2i sin
t for 0 ≤ t ≤ 2. Now  Re
z
t  e  = e2 cos
t ≤ e2 for 0 ≤ t ≤ 2. Since the length of $ is 4, then 
  Re
z  dz ≤ 4e2   e $

This number bounds the magnitude of the integral. It is not claimed to be an approximation to the value of the integral to any degree of accuracy.

22.2.3 Integrals of Series of Functions We often want to interchange an integral and a series. We would like conditions under which  

?  fn
zdz = fn
zdz $

n=1

n=1 $

We willshow that, if we can bound each fn
z, for z on the curve, by a positive constant Mn so that n=1 Mn converges, then we can interchange the summation and the integral and integrate the series term-by-term. THEOREM 22.6

Term-by-Term Integration

Let $ be a smooth curve and let fn be continuous on $ for n = 1 2     Suppose for each positive integer n there is a positive number Mn such that n=1 Mn converges and, for all z on $, fn
z ≤ Mn 

22.2 The Integral of a Complex Function Then then



n=1 fn
z

converges absolutely for all z on $. Further, if we denote
$

g
zdz =


 n=1 $



n=1 fn
z

989

= g
z,

fn
zdz

 For each  z on $, the real series n=1 fn
z converges by comparison with the convergent series n=1 Mn . Now let L be the length of $ and consider the partial sum

Proof

FN
z =

N 

fn
z

n=1

Each FN is continuous on $ and   
 
 N
     fn
zdz =  g
zdz − FN
zdz  g
zdz −  $  $ $ n=1   ≤ L max g
z − FN
z  z on $

Now for all z on $,

              fn
z ≤  Mn   g
z − Fn
z =  n=N +1  n=N +1   If   is any positive number, we can choose N large enough that n=N +1 Mn < /L, because n=1 Mn converges. But then  max g
z − FN
z <  z on $ L so   
 N
    fn
zdz < L =   g
zdz −  $  L n=1 $ for N sufficiently large. This proves that
N
 lim fn
zdz = g
zdz N →

n=1 $

$

as we wanted to show. Of course, the theorem applies to a power series within its circle of convergence, with fn
z = cn
z − z0 n .

SECTION 22.2

PROBLEMS

 In each of Problems 1 through 15, evaluate $ f
zdz. All closed curves are oriented counterclockwise unless specific exception is made. 1. f
z = 1 $
t = t2 − it for 1 ≤ t ≤ 3 2. f
z = z2 − iz $ is the quarter circle about the origin from 2 to 2i

3. f
z = Re
z $ is the line segment from 1 to 2 + i 4. f
z = 1/z $ is the part of the half-circle of radius 4 about the origin, from 4i to −4i 5. f
z = z − 1 $ is any piecewise smooth curve from 2i to 1 − 4i 6. f
z = iz2  $ is the line segment from 1 + 2i to 3 + i

990

CHAPTER 22

Complex Integration

7. f
z = sin
2z $ is the line segment from −i to −4i

13. f
z = iz $ is the line segment from 0 to −4 + 3i

8. f
z = 1 + z  $ is the part of the circle of radius 3 about the origin from −3i to 3i

14. f
z = Im
z $ is the circle of radius 1 about the origin

9. f
z = −i cos
z $ is any piecewise smooth curve from 0 to 2 + i

15. f
z = z2  $ is the line segment from −i to 1   16. Find a bound for  $ cos
z2 dz, if $ is the circle of radius 4 about the origin.    1  dz, if $ is the line segment 17. Find a bound for  $ 1+z from 2 + i to 4 + 2i.

2

10. f
z = z2  $ is the line segment from −4 to i 11. f
z =
z − i3  $
t = t − it2 for 0 ≤ t ≤ 2 12. f
z = eiz  $ is any piecewise smooth curve from −2 to −4 − i

22.3

Cauchy’s Theorem Cauchy’s theorem (or the Cauchy integral theorem) is considered the fundamental theorem of complex integration, and is named for the early nineteenth-century French mathematician and engineer Augustin-Louis Cauchy. He had the idea of the theorem, as well as many of its consequences, but was able to prove it only under what were later found to be unnecessarily restrictive conditions. Edouard Goursat proved the theorem as it is now usually stated, and for this reason it is sometimes called the Cauchy–Goursat theorem. The statement of the theorem implicitly makes use of the Jordan curve theorem, discussed previously in Section 13.2 in connection with Green’s theorem. It states that a continuous simple closed curve $ in the plane separates the plane into two open sets. One of these sets is unbounded, and is called the exterior of $, and the other set is bounded and is called the interior of $. The (graph of the) curve itself does not belong to either of these sets, but forms the boundary for both. Figure 22.16 illustrates the theorem. Although this conclusion may seem obvious for closed curves we might routinely sketch, it is difficult to prove because of the generality of its statement. Some terminology will make the statement of Cauchy’s theorem more efficient.

DEFINITION 22.4

Path

A path is a simple, piecewise smooth curve. A path in a set S is a path whose graph lies in S.

Thus, a path is a join of smooth curves with no self-intersections.

y

Exterior of 

 Interior of 

x FIGURE 22.16

theorem.

Jordan curve

22.3 Cauchy’s Theorem

DEFINITION 22.5

991

Connected Set

A set S of complex numbers is connected if, given any two points z and w in S, there is a path in S having z and w as end points.

S is connected if it is possible to get from any point of S to any other point by moving along some path lying completely in S. An open disk is connected, as is a closed disk, while the set consisting of the two open disks z < 1 and z − 10i < 1 is not (Figure 22.17), because we cannot get from 0 to 10i without going outside the set.

DEFINITION 22.6

Domain

An open, connected set of complex numbers is called a domain.

D is a domain if 1. about any z in D, there is some open disk containing only points of D, and 2. we can get from any point in D to any other point in D by a path in D. For example, any open disk is a domain, as is the upper half plane consisting of all z with Im
z > 0. A closed disk is not a domain (it is connected but not open), and a set consisting of two disjoint open disks is not a domain (it is open but not connected).

DEFINITION 22.7

Simply Connected

A set S of complex numbers is simply connected if every closed path in S encloses only points of S.

Every open disk is simply connected (Figure 22.18). If we draw a closed path in an open disk, this closed path will enclose only points in the open disk. The annulus of Figure 22.19, consisting of points between two concentric circles, is not simply connected, even though it is connected. We can draw a closed path contained in the annulus, but enclosing the inner boundary circle of the annulus. This curve encloses points not in the annulus, namely those enclosed by the inner boundary circle. We are now ready to state one version of Cauchy’s theorem. THEOREM 22.7

Cauchy’s Theorem

Let f be differentiable on a simply connected domain G. Let $ be a closed path in G. Then
f
zdz = 0 $ / Often integrals around closed paths are denoted . In this notation, the conclusion of the / theorem reads $ f
zdz = 0. The oval on the integral sign is just a reminder that the path is closed, and does not alter the way we operate with the integral or the way it is evaluated.

992

CHAPTER 22

Complex Integration y

10i y z  1

y

x x

x

FIGURE 22.17

FIGURE 22.18

Disjoint open disks form a set that is not connected.

An open disk is simply connected.



z  2

FIGURE 22.19 The set of points between two concentric circles is not simply connected.

/ Informally, Cauchy’s theorem states that $ f
zdz = 0 if f is differentiable on the curve and on all points enclosed by the curve. As a convention, closed curves are understood to be oriented positively (counterclockwise), unless specific exception is made.

EXAMPLE 22.15

/ 2 Evaluate $ ez dz, where $ is any closed path in the plane. 2 Figure 22.20 shows a typical $. Here f
z = ez is differentiable for all z, and the entire plane is a simply connected domain. Therefore - 2 ez dz = 0 $

EXAMPLE 22.16

Evaluate

$

2z + 1 dz z2 + 3iz

where $ is the circle z + 3i = 2 of radius 2 and center −3i (Figure 22.21). y

y 

x  x

2 z  3i  2

FIGURE 22.20 A simple closed path $.

FIGURE 22.21

22.3 Cauchy’s Theorem

993

We can parametrize $
t = −3i + 2eit for 0 ≤ t ≤ 2. $
t traverses the circle once counterclockwise as t varies from 0 to 2. First observe that f
z is differentiable except at points where the denominator vanishes, 0 and −3i. Use a partial fractions decomposition to write   11 6+i 1 f
z = +  3i z 3 z + 3i Since 1/z is differentiable on and within the simply connected domain enclosed by $, by Cauchy’s theorem, - 11 dz = 0 $ 3i z However, 1/
z + 3i is not differentiable in the simply connected domain enclosed by $, so Cauchy’s theorem does not apply to an integral of this function. We will evaluate this integral directly by writing z
t = −3i + 2eit : - 6+i 1 6 + i
2 1 dz = z

tdt 3 z + 3i 3 z
t + 3i $ 0 6 + i
2 1 6 + i
2 6+i
2i = 2ieit dt = idt = it 3 0 2e 3 0 3 Therefore $

  6+i 2 2z + 1 dz =
2i = − + 4i  z2 + 3iz 3 3

We will see more impressive ramifications of Cauchy’s theorem shortly.

22.3.1

Proof of Cauchy’s Theorem for a Special Case

If we add an additional hypothesis, Cauchy’s theorem is easy to prove. Let f
z = u
x y + iv
x y and assume that u and v and their first partial derivatives are continuous on G. Now we obtain Cauchy’s theorem immediately by applying Green’s theorem and the Cauchy–Riemann equations to equation (22.1). If D consists of all points on and enclosed by $, then f
zdz = udx − vdy + i vdx + udy $

$

$



 u v 


−v u  − dA + i − dA = 0 =

x

y

x y D D because, by the Cauchy–Riemann equations,

u v

u

v = and =− 

x y

y

x This argument is good enough for many settings in which Cauchy’s theorem is used. However, it is not an optimal argument because it makes the additional assumption about continuity of the partial derivatives of u and v. A rigorous proof of the theorem as stated involves topological subtleties we do not wish to engage here. In the next section we will develop some important consequences of Cauchy’s theorem.

CHAPTER 22

994

Complex Integration

SECTION 22.3

PROBLEMS 5. f
z = z $ is the unit circle about the origin

In each of Problems 1 through 12, evaluate the integral of the function over the given closed path. All paths are positively oriented (counterclockwise). In some cases Cauchy’s theorem applies, and in some it does not.

7. f
z = zez  $ is the circle z − 3i = 8 8. f
z = z2 − 4z + i $ is the rectangle with vertices 1, 8, 8 + 4i and 1 + 4i

1. f
z = sin
3z $ is the circle z = 4 2z  $ is the circle z − i = 3 2. f
z = z−i 1 3. f
z =  $ is given by z − 2i = 2
z − 2i3

9. f
z = z2  $ is the circle of radius 7 about the origin 10. f
z = sin
1/z $ is the circle z − 1 + 2i = 1 11. f
z = Re
z $ is given by z = 2

4. f
z = z sin
z $ is the square having vertices 0, 1, 1 + i and i 2

22.4

6. f
z = 1/z  $ is the circle of radius 5 about the origin

12. f
z = z2 + Im
z $ is the square with vertices 0, −2i, 2 − 2i and 2

Consequences of Cauchy’s Theorem This section lays out some of the main results of complex integration, with profound implications for understanding the behavior and properties of complex functions, as well as for applications of the integral. As usual, all integrals over closed curves are taken with a counterclockwise orientation, unless otherwise noted.

22.4.1 Independence of Path In Section 22.2.2 we mentioned  independence of path, according to which, under certain conditions on f , the value of $ f
zdz depends only on the end points of the curve, and not on the particular curve chosen between those end points. Independence of path can also be viewed from the perspective of Cauchy’s theorem. Suppose f is differentiable on a simply connected domain G, and z0 and z1 are points of G. Let $1 and $2 be piecewise smooth curves in G having initial point z0 and terminal point z1 (Figure 22.22). If we reverse the orientation on $2 , we obtain the new curve, −$2 , going from z1 to z0 . Further, the join of $1 and −$2 forms a closed curve $, having initial and terminal point z0 (Figure 22.23). By Cauchy’s theorem and Theorem 22.3, $

f
zdz = 0 =

$1 ⊕
−$2 

f
zdz =

$1

f
zdz −

f
zdz $2

implying that $1

f
zdz =

f
zdz $2

This means that the integral does not depend on the particular curve (in G) between z0 and z1 , and is therefore independent of path. This argument is not entirely rigorous, because $1 ⊕
−$2  need not be a simple curve (Figure 22.24). In fact, $1 and $2 may cross each other any number of times as they progress from z0 to z1 . Nevertheless, we wanted to point out the connection between Cauchy’s theorem and the concept of independence of path of an integral which was discussed previously.

22.4 Consequences of Cauchy’s Theorem 1 z1

z1

2


2

1

If write

 $

z2

2

1 z0

995

z0

z0

FIGURE 22.22

FIGURE 22.23

FIGURE 22.24

Paths $1 and $2 from z0 to z1 .

Closed curve $ = $1 ⊕
−$2 .

$1 ⊕
−$2  need not be simple.

f
zdz is independent of path in G, and $ is any path from z0 to z1 , we sometimes $

f
zdz =




z1

f
zdz z0

The symbol on the right has the value of the line integral on the left, with $ any path from z0 to z1 in G.

22.4.2

The Deformation Theorem

The deformation theorem enables us, under certain conditions, to replace one closed path of integration with another, perhaps more convenient one. THEOREM 22.8

Deformation Theorem

Let $ and ! be closed paths in the plane, with ! in the interior of $. Let f be differentiable on an open set containing both paths and all points between them. Then,

f
zdz = f
zdz $

!

Figure 22.25 shows the setting of the theorem. We may think of deforming one curve, say !, to the other. Imagine ! as made of rubber, and continuously deform it into the shape of $. In doing this, it is necessary that the intermediate stages of the deformation from ! to $ only





FIGURE 22.25 Deforming ! continuously into $.

996

CHAPTER 22

Complex Integration

cross over points at which f is differentiable, hence the hypothesis about f being differentiable on some open set containing both paths and all points between them. The theorem states that the integral of f has the same value over both paths when one can be deformed into the other, moving only over points at which the function is differentiable. This means that we can replace $ with another path ! that may be more convenient to use in evaluating the integral. Consider the following example.

EXAMPLE 22.17

Evaluate

$

1 dz z−a

over any closed path enclosing the given complex number a. Figure 22.26 shows a typical such path. We cannot parametrize $ because we do not know it specifically – it is simply any path enclosing a. Let ! be a circle of radius r about a, with r small enough that ! is enclosed by $ (Figure 22.27). Now f
z = 1/
z − a is differentiable at all points except a, hence on both curves and the region between them. By the deformation theorem, - 1 - 1 dz = dz $ z−a ! z−a But ! is easily parametrized: !
t = a + reit for 0 ≤ t ≤ 2. Then - 1
2 1
2 it dz = ire dt = idt = 2i reit ! z−a 0 0 Therefore

-

1 dz = 2i z−a

$

The point is that, by means of the deformation theorem, we can evaluate this integral over any path enclosing a. Of course, if $ does not/ enclose a, and a is not on $, then 1/
z − a is differentiable on $ and the set it encloses, so $ 1/
z − adz = 0 by Cauchy’s theorem. The proof of the theorem employs a technique we will find useful in several settings. Figure 22.28 shows graphs of typical paths $ and !. Insert lines L1 and L2 between $ and ! (Figure 22.29) and use these to form two closed paths # and + (shown separated for emphasis in Figure 22.30). One path, #, consists of parts of $ and !, together with L1 and L2 , with orientation on each piece as shown in order to have positive orientation on #. The other path, + , consists of the rest of $ and !, again with L1 and L2 , with the orientation chosen on each piece so that + has positive orientation. Figure 22.31 shows the paths more realistically, Proof

y y

 a

 a

x

x FIGURE 22.26



FIGURE 22.27

22.4 Consequences of Cauchy’s Theorem

997

y 



 x

FIGURE 22.28



y

L1

y L1 L1



L2

L2 L2

 L1

L2 x



x

FIGURE 22.29

FIGURE 22.30

 FIGURE 22.31

sharing the inserted segments L1 and L2 . In Figure 22.31, $ is oriented counterclockwise, but ! is clockwise, due to their orientations as parts of # and + . Because f is differentiable on both # and + and the sets they enclose, Cauchy’s theorem yields f
zdz = f
zdz = 0 #

Then

#

+

f
zdz +

+

f
zdz = 0

(22.2)

In this sum of integrals, each L1 and L2 is integrated over in one direction as part of # and the opposite direction as part of + . The contributions from these segments therefore cancel in the sum (22.2). Next observe that, in adding these integrals, we obtain the integral over all of $, oriented counterclockwise, together with the integral over all of !, oriented clockwise. In view of Theorem 22.3, equation (22.2) becomes f
zdz − f
zdz = 0 #

or

!

#

f
zdz =

f
zdz !

in which the orientation on both $ and ! in these integrals is positive (counterclockwise). This proves the theorem.

22.4.3

Cauchy’s Integral Formula

We will now state a remarkable result which gives an integral formula for the values of a differentiable function. THEOREM 22.9

Cauchy Integral Formula

Let f be differentiable on an open set G. Let $ be a closed path in G enclosing only points of G. Then, for any z0 enclosed by $, 1 - f
z f
z0  = dz 2i $ z − z0 We will see many uses for this theorem, but one is immediate. Write the formula as - f
z dz = 2if
z0  $ z − z0

998

CHAPTER 22

Complex Integration

This gives, under the conditions of the theorem, an evaluation of the integral on the left as a constant multiple of the function value on the right.

EXAMPLE 22.18

Evaluate -

2

ez dz $ z−i for any closed path that does not pass through i. 2 Let f
z = ez . Then f is differentiable for all z. There are two cases. / e z2 2 Case 1 $ does not enclose i. Now $ z−i dz = 0 by Cauchy’s theorem, because ez /
z − 2 is differentiable on and within $. Case 2 $ encloses i. By Cauchy’s integral formula, with z0 = i, - e z2 dz = 2if
i = 2ie−1  $ z−i

EXAMPLE 22.19

Evaluate

- e2z sin
z2  dz z−2 $

over any path not passing through 2. Let f
z = e2z sin
z2 . Then f is differentiable for all z. This leads to two cases. Case 1 If $ does not enclose 2, then f
z/
z − 2 is differentiable on the curve and at all points it encloses. Now the integral is zero by Cauchy’s theorem. Case 2 If $ encloses 2, then by the integral formula, - e2z sin
z2  dz = 2if
2 = 2ie4 sin
4 z−2 $ Observe the distinction between the roles of f
z in Cauchy’s theorem and in Cauchy’s / f
zdz. The integral repreintegral representation. Cauchy’s theorem is concerned with $ / sentation is concerned with integrals of the form $ f
z/
z − z0 dz, with f
z given, but multiplied by a factor 1/
z − z0  which is not defined at z0 . If $ does not enclose z0 , then f
z/
z − z0  = g
z may be differentiable at z0 and we can attempt to apply Cauchy’s theorem / If z0 is enclosed by $, then under the appropriate conditions, the integral formula to $ g
zdz. / gives $ g
zdz in terms of f
z0 . Here is a proof of the integral representation. First use the deformation theorem to replace $ with a circle ! of radius r about z0 , as in Figure 22.32. Then - f
z - f
z − f
z  + f
z  - f
z 0 0 dz = dz = dz z − z0 $ z − z0 ! z − z0 ! - f
z − f
z  1 0 dz + dz = f
z0  z − z0 ! z − z0 !

Proof

22.4 Consequences of Cauchy’s Theorem

999

y !

 x

z0

FIGURE 22.32

in which f
z0  could be brought outside the first integral because f
z0  is constant. By Example 22.17, 1 dz = 2i ! z − z0 because ! encloses z0 . Therefore - f
z $

z − z0

dz = 2if
z0  +

- f
z − f
z  0 dz z − z0 !

The integral representation is proved if we can show that the last integral is zero. Write !
t = z0 + reit for 0 ≤ t ≤ 2. Then     f
z − f
z0   
2 f
z0 + reit  − f
z0  it     dz =  ire dt  ! z − z0 reit 0 
2    it  =

f
z0 + re  − f
z0 dt 0

≤




2

0

≤ 2

  f
z0 + reit  − f
z0  dt



   it   max f
z0 + re  − f
z0  

0≤t≤2

By continuity of f
z at z0 , f
z0 + reit  → f
z0  as r → 0, so the term on the right in this inequality has limit zero as r → 0 Therefore we can make   f
z − f
z0    dz  ! z−z 0

arbitrarily small by choosing r sufficiently small. But this integral is independent of r by the deformation theorem. Hence   f
z − f
z0    dz = 0  ! z − z0 so - f
z − f
z  0 dz = 0 z − z0 ! and the theorem is proved. The integral representation gives some idea of how strong the condition of differentiability is for complex functions. The integral - f
z dz $ z − z0 equals 2if
z0 , and so determines f
z0  at each z0 enclosed by $. But the value of this integral depends only on the values of f
z on $. Thus, for a differentiable function, knowing function

1000

CHAPTER 22

Complex Integration

values on $ determines the values of the function at all points enclosed by $. There is no analogous result for differentiable real functions. Knowing the values of a differentiable real function 
x at the end points of an interval in general gives no information about values of this function throughout the interval.

22.4.4 Cauchy’s Integral Formula for Higher Derivatives We will now show that a complex function that is differentiable on an open set must have derivatives of all orders on this set. There is no result like this for real functions. A real function that is differentiable need not have a second derivative. And if it has a second derivative, it need not have a third, and so on. Not only does a differentiable complex function have derivatives of all orders, we will show that the nth derivative of the function at a point is also given by an integral formula, very much in the spirit of Cauchy’s integral formula. THEOREM 22.10

Cauchy’s Integral Formula for Higher Derivatives

Let f be differentiable on an open set G. Then f has derivatives of all orders at each point of G. Further, if $ is a closed path in G enclosing only points of G, and z0 is any point enclosed by $, then f
z n! f
n
z0  = dz 2i $
z − z0 n+1 The integral on the right is exactly what we would obtain by differentiating Cauchy’s integral formula for f
z0 , n times with respect to z0 , under the integral sign. As with the integral formula, this conclusion is often used to evaluate integrals

EXAMPLE 22.20

Evaluate -

3

$

ez dz
z − i3

with $ any path not passing through i. If $ does not enclose i then this integral is zero by Cauchy’s theorem, since the only point 3 at which ez /
z − i3 fails to be differentiable is i. Thus suppose $ encloses i. Because the factor 3 z − i occurs to the third power in the denominator, use n = 2 in the theorem, with f
z = ez , to get 3 ez 2i
2 f
i = if


i dz = 3
z − i 2! $ Now f

z = 3z2 ez and f


z = 6zez + 9z4 ez  3

3

3

so $



ez dz = i 6ie−i + 9e−i =
−6 + 9ie−i  3
z − i 3

The theorem can be proved by induction on n, but we will not carry out the details.

22.4 Consequences of Cauchy’s Theorem

22.4.5

1001

Bounds on Derivatives and Liouville’s Theorem

Cauchy’s integral formula for higher derivatives can be used to obtain bounds on derivatives of all orders. THEOREM 22.11

Let f be differentiable on an open set G. Let z0 be a point of G and let the open disk of radius r about z0 be in G. Suppose f
z ≤ M for z on the circle of radius r about z0 . Then, for any positive integer n,  
n f
z0  ≤ Mn!  rn Let !
t = z0 + reit for 0 ≤ t ≤ 2. Then f
z0 + reit  ≤ M for 0 ≤ t ≤ 2. By Theorems 22.10 and 22.5,      
n  f
z n! 
2 f
z0 + reit  it   f
z0  = n!  dz = ire dt 2  !
z − z n+1  2  0 r n+1 ei
n+1t Proof

0

≤

n!
2 f
z0 + reit  1 n! Mn!
2M n = n  dt ≤ 2 0 rn 2 r r

As an application of this theorem, we will prove Liouville’s theorem on bounded, differentiable functions. THEOREM 22.12

Liouville

Let f be a bounded function that is differentiable for all z. Then f is a constant function. Previously we noted that sin
z is not a bounded function in the complex plane, as it is on the real line. This is consistent with Liouville’s theorem. Since sin
z is differentiable for all z, and is clearly not a constant function, it cannot be bounded. Here is a proof of Liouville’s theorem. Suppose f
z ≤ M for all complex z. Choose any number z0 and any r > 0. By Theorem 22.11, with n = 1,

Proof

M  r Since f is differentiable in the entire complex plane, r may be chosen as large as we like, so f

z0  must be less than any positive number. We conclude that f

z0  = 0, hence f

z0  = 0. Since z0 is any number, then f
z = constant. f

z0  ≤

Liouville’s theorem can be used to give a simple proof of the fundamental theorem of algebra. This theorem states that any nonconstant complex polynomial p
z = a0 + a1 z + · · · + an zn has a complex root. That is, for some number ", p
" = 0. From this it can be further shown that, if an = 0, then p
z must have exactly n roots, counting each root k times in the list if its multiplicity is k. For example, p
z = z2 − 6z + 9 has exactly two roots, 3 and 3 (a root of multiplicity 2). This fundamental theorem assumes only elementary terminology for its statement, and is usually included, in some form, in the high school mathematics curriculum. The leading nineteenth century mathematician Carl Friedrich Gauss considered the theorem so important

1002

CHAPTER 22

Complex Integration

that he devised many proofs (some say nearly twenty) over his lifetime. But even today rigorous proofs of the theorem require mathematical terms and devices far beyond those needed to state it. To prove the theorem using Liouville’s theorem, suppose p
z is a nonconstant complex polynomial and that p
z = 0 for all z. Then 1/p
z is differentiable for all z. Let p
z = a0 + a1 z + · · · + an zn with n ≥ 1 and an = 0. We will show that p
z is bounded for all z. Since an zn = p
z − a0 − a1 z − · · · − an−1 zn−1  then an  zn ≤ p
z + a0  + a1  z + · · · + an−1  zn−1  Then, for z ≥ 1,

 p
z ≥ an  zn − a0  + a1  z + · · · + an−1  zn−1   a1  a  a0  = zn−1 an  z − n−1 − n−2 − · · · − n−1 zn−n z z ≥ zn−1
an  z − a0  − a1  − · · · − an−1  

But then 1 1 ≤ n−1 →0 p
z z
an  z − a0  − a1  − · · · − an−1  as z → . Therefore limn→ 1/ p
z = 0. This implies that, for some positive number R, 1 < 1 if z > R p
z But the closed disk z ≤ R is compact, and 1/ p
z is continuous, so 1/ p
z is bounded on this disk by Theorem 21.1. Therefore, for some M, 1 ≤ M for z ≤ R p
z Now we have 1/ p
z bounded both inside and outside z ≤ 1. Putting these bounds together, 1 ≤ M + 1 for all z p
z both in z ≤ R and in z ≥ R. This makes 1/p
z a bounded function that is differentiable for all z. By Liouville’s theorem, 1/p
z must be constant, a contradiction. Therefore there must be some complex " such that p
" = 0, proving the fundamental theorem of algebra. Complex analysis provides several proofs of this theorem. Later we will see one using a technique for evaluating real integrals of rational functions involving sines and cosines.

22.4.6 An Extended Deformation Theorem The deformation theorem allows us to deform one closed path $ of integration into another, !, without changing the value of the line integral of a differentiable function f . A crucial condition for this process is that no stage of the deformation should pass over a point at which f fails to be differentiable. This means that f needs to be differentiable on both curves and the region between them. We will now extend this result to the case that $ encloses any finite number of disjoint closed paths. As usual, unless explicitly stated otherwise, all closed paths are assumed to be oriented positively (counterclockwise).

22.4 Consequences of Cauchy’s Theorem THEOREM 22.13

1003

Extended Deformation Theorem

Let $ be a closed path. Let !1      !n be closed paths enclosed by $. Assume that no two of $ !1      !n intersect, and no point interior to any !j is interior to any other !k . Let f be differentiable on an open set containing $, each !j , and all points that are both interior to $ and exterior to each !j . Then, $

f
zdz =

n 

f
zdz

j=1 !j

This is the deformation theorem in the case n = 1. Figure 22.33 shows a typical scenario covered by this theorem. With the curves as shown (and the differentiability assumption on f ), the integral of f about $ is the sum of the integrals of f about each of the closed curves !1      !n . We will outline a proof after an illustration of a typical use of the theorem.

EXAMPLE 22.21

Consider

$

z dz
z + 2
z − 4i

where $ is a closed path enclosing −2 and 4i. We will evaluate this integral using the extended deformation theorem. Place a circle !1 about −2 and a circle !2 about 4i, of sufficiently small radii that neither circle intersects the other or $ and each is enclosed by $ (Figure 22.34). Then z z z dz = dz + dz
z + 2
z − 4i
z + 2
z − 4i
z + 2
z − 4i $ !1 !2 Use a partial fractions decomposition to write 1 − 2i 4 + 2i z = 5 5 +5 5 
z + 2
z − 4i z+2 z − 4i

y

4i



!2 y

 !n

!1


2

!2

!1

!3 x FIGURE 22.33

FIGURE 22.34

x

1004

CHAPTER 22

Complex Integration

Then $

z dz =
z + 2
z − 4i



   1 2 4 2 1 1 − i dz + + i dz 5 5 5 5 !1 z + 2 !1 z − 4i     1 2 4 2 1 1 + − i dz + + i dz 5 5 z + 2 5 5 z − 4i !2 !2

On the right, the second and third integrals are zero by Cauchy’s theorem (!1 does not enclose 4i and !2 does not enclose −2). The first and fourth integrals equal 2i by Example 22.17. Therefore     1 2 4 2 z dz = 2i − i + + i = 2i 5 5 5 5 $
z + 2
z − 4i

A proof of the theorem can be modeled after the proof of the deformation theorem. As suggested in Figure 22.35, draw line segments L1 from $ to !1 , L2 from !1 to !2      Ln from !n−1 to !n and, finally, Ln+1 from !n to $. Form the closed paths # and shown separately in Figures 22.36 and 22.37, and as they actually appear in Figure 22.38. Then f
zdz + f
zdz = 0

Proof

#



both integrals being zero by Cauchy’s theorem. (By the hypotheses of the theorem, f is differentiable on and inside both # and ). In this sum of integrals over # and , each line segment Lj is integrated over in both directions, hence the contributions of the integrals over these segments is zero. Further, in this sum we retrieve the integral of f
z over all of $, and all of each !j , with the orientation counterclockwise on $ and clockwise on each !j (note the orientations in Figure 22.38). y L n1

!1

!n

!2

Ln

L1



L2

!n
1

x

L n
1

 FIGURE 22.35

FIGURE 22.36

 FIGURE 22.37

FIGURE 22.38

22.4 Consequences of Cauchy’s Theorem

1005

Reversing the orientations on the !j
s, so that all paths are oriented counterclockwise, the last sum becomes n  f
zdz − f
zdz = 0 $

j=1 !j

yielding the conclusion of the theorem.

SECTION 22.4

PROBLEMS

 In each of Problems 1 through 12, evaluate f
zdz for the given function and closed (positively oriented) path. These problems may involve Cauchy’s theorem, Cauchy’s integral formulas and/or the deformation theorems. z4  $ is any closed path enclosing 2i 1. f
z = z − 2i 2. f
z =

sin
z2   $ is any closed path enclosing 5 z−5

3. f
z =

z2 − 5z + i  $ is the circle z = 3 z − 1 + 2i

2z3  $ is the rectangle having vertices 4 ± i
z − 22 and −4 ± i

4. f
z =

5. f
z =

iez  $ is the circle z − 1 = 4
z − 2 + i2

6. f
z =

cos
z − i  $ is any closed path enclosing −2i
z + 2i3

7. f
z =

z sin
3z  $ is the circle z − 2i = 9
z + 43

8. f
z = 2iz z  $ is the line segment from 1 to −i

9. f
z = − ing −4


2 + i sin
z4   $ is any closed path enclos
z + 42

10. f
z =
z − i2  $ is the semicircle of radius 1 about 0 from i to −i 11. f
z = Re
z + 4 $ is the line segment from 3 + i to 2 − 5i 12. f
z =

3z2 cosh
z  $ is the circle of radius 8 about 1
z + 2i2

13. Evaluate




2

ecos
 cos
sin
d

0

/ Hint: Consider $
ez /zdz, with $ the unit circle about the origin. Evaluate this integral once using Cauchy’s integral formula, then again directly by using the coordinate functions for $. 14. Use the extended form of the deformation theorem / z − 4i to evaluate $ 3 dz, where $ is a closed path z + 4z enclosing the origin, 2i and −2i.

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CHAPTER

23

POWER SERIES REPRESENTATIONS ISOLATED ZEROS AND THE IDENTITY THEOREM THE MAXIMUM MODULUS THEOREM THE LAURENT EXPANSION POWER SERIES REPRESENTATIONS ISOLATED

Series Representations of Functions

We will now develop two kinds of representations of a f
z in series of powers of z − z0 . The first series will contain only nonnegative integer powers, hence is a power series, and applies when f is differentiable at z0 . The second will contain negative integer powers of z − z0 as well, and will be used when f is not differentiable at z0 .

23.1

Power Series Representations We know that a power series that converges in an open disk, or perhaps the entire plane, defines a function that is infinitely differentiable within this disk or the plane. We will now go the other way and show that a function that is differentiable on an open disk is represented by a power series expansion about the center of this disk. This will have important applications, including information about zeros of functions, and the maximum value that can be taken by the modulus f
z of a differentiable function.

THEOREM 23.1

Taylor Series

Let f be differentiable on an open disk D about z0 . Then, for each z in D, f
z =

 f
n
z0  n=0

n!


z − z0 n 

The series on the right is the Taylor series of f
z about z0 , and the number f
n
z0 /n! is the nth Taylor coefficient of f
z at z0 . The theorem asserts that the Taylor series of f
z converges to f
z, hence represents f
z, within this disk. 1007

CHAPTER 23

1008

Series Representations of Functions

Let z be any point of D and R the radius of D. Choose a number r with z − z0  < r < R and let ! be the circle of radius r about z0 (Figure 23.1). By the Cauchy integral formula, using w for the variable of integration on !, Proof

R

r

!

z0 w

f
z = FIGURE 23.1

1 - f
w dw 2i ! w − z

Now 1 1 1 1 = =  w − z w − z0 −
z − z0  w − z0 1 −
z − z0 /
w − z0  Since w is on !, and z is enclosed by !, then    z − z0    w−z  < 1 0 so we can use a convergent geometric series to write   1 z − z0 n  1  1 = =
z − z0 n  n+1 w − z w − z0 n=0 w − z0
w − z  0 n=0 Then  f
w f
w =
z − z0 n  w − z n=0
w − z0 n+1

(23.1)

Since f is continuous on !, for some M, f
w ≤ M for w on !. Further, w − z0  = r, so       f
w 1 z − z0  n n   
w − z n+1
z − z0   ≤ M r r 0 Designate M

1 r



z − z0  r

n = Mn 

 Then n=0 Mn converges (this series is a constant times a convergent geometric series). By Theorem 22.6, the series in equation (23.1) can be integrated term-by-term to yield   1 -  f
w 1 - f
w n dw =
z − z0  dw f
z = 2i ! w − z 2i ! n=0
w − z0 n+1     1 f
w f
n
z0  n =
z − z0 n  dw
z − z  = 0 n+1 2i
w − z  n! ! 0 n=0 n=0 in which we used Cauchy’s integral formula for the nth derivative to write the coefficient in the last series. This proves the theorem. A complex function is said to be analytic at z0 if it has a power series expansion in some open disk about z0 . We have just proved that a function that is differentiable on an open disk about a point is analytic at that point. We only compute the coefficients of a Taylor series by derivative or integral formulas when other means fail. When possible, we use known series and operations such as differentiation

23.1 Power Series Representations

1009

and integration to derive a series representation. This strategy makes use of the uniqueness of power series representations. THEOREM 23.2

Suppose, in some disk z − z0  < r, 

cn
z − z0 n =

n=0



dn
z − z0 n 

n=0

Then, for n = 0 1 2     cn = dn . Proof

If we let f
z be the function defined in this disk by both power series, then cn =

f
n
z0  = dn  n!

This means that, no matter what method is used to find a power series for f
z about z0 , the end result is the Taylor series.

EXAMPLE 23.1

Find the Taylor expansion of ez about i. We know that, for all z, ez =

 1 n z  n=0 n!

For an expansion about i, the power series must be in terms of powers of z − i. Thus write ez = ez−i+i = ei ez−i =



ei

n=0

1
z − in  n!

This series converges for all z. In this example, it would have been just as easy to compute the Taylor coefficients directly: cn =

ei f
n
i =  n! n!

EXAMPLE 23.2

Write the Maclaurin series for cos
z3 . A Maclaurin expansion is a Taylor series about zero. For all z, cos
z =


−1n n=0


2n!

z2n 

3

All we need to do is replace z with z : cos
z3  =


−1n n=0


2n!


z3 2n =


−1n n=0


2n!

z6n 

Since this is an expansion about the origin, it is the expansion we seek.

1010

CHAPTER 23

Series Representations of Functions

EXAMPLE 23.3

Expand 2i 4 + iz in a Taylor series about −3i. We want a series in powers of z + 3i. Do some algebraic manipulation and then use a geometric series. In order to get z + 3i into the picture, write 2i 2i 2i 1 2i = = =  i 4 + iz 4 + i
z + 3i + 3 7 + i
z + 3i 7 1 + 7
z + 3i If t < 1, then   1 1 = =
−tn =
−1n  1 + t 1 −
−t n=0 n=0

With t =
z + 3ii/7, we have

n   1 i n
z + 3i =
−1  7 1 + 7i
z + 3i n=0

Therefore

 n  2i i 2i  2
−1n in+1 =
−1n
z + 3in =
z + 3in  n+1 4 + iz 7 n=0 7 7 n=0

Because this is a series expansion of the function about −3i, it is the Taylor series about −3i. This series converges for   i  
z + 3i < 1 7  or z + 3i < 7 Thus z must be in the open disk of radius 7 about −3i. The radius of convergence of this series is 7. From Section 21.2, we can differentiate a Taylor series term by term within its open disk of convergence. This is sometimes useful in deriving the Taylor expansion of a function.

EXAMPLE 23.4

Find the Taylor expansion of f
z = 1/
1 − z3 about the origin. We could do this by algebraic manipulation, but it is easier to begin with the familiar geometric series g
z =

 1 = zn 1 − z n=0

for z < 1

Then g

z =

 1 = nzn−1 2
1 − z n=1

23.1 Power Series Representations

1011

and g


z =

  2 n−2 = n
n − 1z =
n + 1
n + 2zn
1 − z3 n=2 n=0

for z < 1. Then f
z =

 1 n=0 2


n + 1
n + 2zn

for z < 1

When f
z is expanded in a power series about z0 , the radius of convergence of the series will be the distance from z0 to the nearest point at which f
z is not differentiable. Think of a disk expanding uniformly from z0 , free to continue its expansion until it hits a point at which f
z is not differentiable. For example, suppose f
z = 2i/
4 + iz, and we want a Taylor expansion about −3i. The only point at which f
z is not defined is 4i, so the radius of convergence of this series will be the distance and 4i, or 7. We obtained this result previously from the Taylor  betweenn −3i n+1 n+1
2
−1 i /7 
z + 3in of f
z. expansion n=0

EXAMPLE 23.5

We will find the radius of convergence of the Taylor series of csc
z about 3 − 4i. Since csc
z = 1/ sin
z, this function is differentiable except at z = n, with n any integer. As Figure 23.2 illustrates,  is the nearest point to 3 − 4i at which csc
z is not differentiable. The distance between  and 3 − 4i is 
 − 32 + 16 and this is the radius of convergence of the expansion of csc
z about 3 − 4i. y 

x


3) 2  16 (w 3
4i FIGURE 23.2

Existence of a power series expansion implies the existence of an antiderivative. THEOREM 23.3

Let f be differentiable in an open disk D about z0 . Then there is a differentiable function F such that F

z = f
z for all z in D. Proof

We know that f has a power series expansion in D: f
z =



cn
z − z0 n 

n=0

Let F
z =



1 cn
z − z0 n+1 n + 1 n=0

1012

CHAPTER 23

Series Representations of Functions

for z in D. It is routine to check that this power series has radius of convergence at least as large as the radius of D, and that F

z = f
z for z in D.

23.1.1 Isolated Zeros and the Identity Theorem We will use the Taylor series representation of a differentiable function to derive important information about its zeros.

DEFINITION 23.1

A number " is a zero of f if f
" = 0. A zero " of f is isolated if there is an open disk about " containing no other zero of f .

For example, the zeros of sin
z are n, with n any integer. These zeros are all isolated. By contrast, let  sin
1/z if z = 0 f
z =  0 if z = 0 The zeros of f are 0 and the numbers 1/n, with n any nonzero integer. However, 0 is not an isolated zero, since any open disk about zero contains other zeros, 1/n, for n sufficiently large. We will show that the behavior of the zeros in this example disqualifies f from being differentiable at 0. THEOREM 23.4

Let f be differentiable on an open disk G and let " be a zero of f in G. Then either " is an isolated zero of f , or there is an open disk about " on which f
z is identically zero. This means that a differentiable function that is not identically zero on some disk can have only isolated zeros. Proof

Write the power series expansion of f
z about ", f
z =



cn
z − "n 

n=0

in some open disk D centered at ". Now consider two cases. Case 1 If each cn = 0, then f
z = 0 throughout D. Case 2 Suppose, instead, that some coefficients in the power series are not zero. Let m be the smallest integer such that cm = 0. Then c0 = c1 = · · · = cm−1 = 0 and cm = 0. For z in D, f
z =

 n=m

cn
z − "n =

 n=0

cn+m
z − "n+m =
z − "m

 n=0

cn+m
z − "n 

23.1 Power Series Representations

1013

 n Now n=0 cn+m
z − " is a power series, and so defines a differentiable function g
z for z in D. Further, f
z =
z − "m g
z But g
" = cm = 0, so there is some open disk K about " in which g
z = 0. Therefore, for z = " in K, f
z = 0 also, hence " is an isolated zero of f . This proof contained additional information about zeros that will be useful later. The number m in the proof is called the order of the zero " of f
z. It is the smallest integer m such that the coefficient cm in the expansion of f
z about " is nonzero. Now recall that cn = f
n
"/n!. Thus c0 = c1 = · · · = cm−1 = 0 implies that f
" = f

" = · · · = f
m−1
" = 0 while cm = 0 implies that f
m
" = 0 In summary, an isolated zero " of f has order m if the function and its first m − 1 derivatives vanish at ", but the mth derivative at " is nonzero. Put another way, the order of the zero " is the lowest order derivative of f that does not vanish at ". We also derived the important fact that, if " is an isolated zero of order m of f , then we can write f
z =
z − "m g
z where g is also differentiable in some disk about ", and g
" = 0.

EXAMPLE 23.6

Consider f
z = z2 cos
z. 0 is an isolated zero of this differentiable function. Compute f

z = 2z cos
z − z2 sin
z and f


z = 2 cos
z − 4z sin
z + z2 cos
z Observe that f
0 = f

0 = 0 while f


0 = 0. Thus 0 is a zero of order 2 of f . In this case, we already have f
z =
z − 02 g
z with g
0 = 0, since we can choose g
z = cos
z

EXAMPLE 23.7

We have to exercise some care in identifying the order of a zero. Consider f
z = z2 sin
z. 0 is an isolated zero of f . Compute f

z = 2z sin
z + z2 cos
z f


z = 2 sin
z + 4z cos
z − z2 sin
z f
3
z = 2 cos
z + 4 cos
z − 4z sin
z − 2z sin
z − z2 cos
z = 6 cos
z − 6z sin
z − z2 cos
z

1014

CHAPTER 23

Series Representations of Functions

Then f
0 = f

0 = f


0 = 0 while f
3
0 = 0 This means that 0 is a zero of order 3. We can write  
−1n 2n+1
−1n 2n f
z = z2 z z = z3 g
z = z3
2n + 1!
2n + 1! n=0 n=0 where g
z =


−1n 2n z n=0
2n + 1!

is differentiable (in this example, for all z) and g
0 = 1 = 0. As a result of Theorem 23.3, we can show that, if a differentiable complex function vanishes on a convergent sequence of points in a domain (open, connected set), then the function is identically zero throughout the entire domain. This is a very strong result, for which there is no analogue for real functions. For example, consider  0 for x ≤ 0  h
x = 2 for x > 0 x whose graph is shown in Figure 23.3 This function is differentiable for all real x, and is identically zero on half the line, but is not identically zero over the entire line. Another difference between differentiability for real and complex functions is evident in this example. While h is differentiable, it does not have a power series expansion about 0. By contrast, a complex function that is differentiable on an open set has a power series expansion about each point of the set. THEOREM 23.5

Let f be differentiable on a domain G. Suppose zn  is a sequence of distinct zeros of f in G, converging to a point of G. Then f
z = 0 for all z in G.

h(x) 16 14 12 10 8 6 4 2
4


2

0

2

FIGURE 23.3



h
x =

x2 0

for x > 0 < 1 . for x ≤ 0

4

x

23.1 Power Series Representations

G D3 y

D2 ξ3 ξ 2 ξ1 "

D1 f(z)  0

Dn D0

f(z)  0

w

"

ξn
1

D x FIGURE 23.4

1015

Distance  2r

Dn
1 FIGURE 23.5

D0 " z1 D1 FIGURE 23.6

Suppose zn → " in G. Since f is continuous, f
zn  → f
". But each f
zn  = 0, so f
" = 0 also, and " must also be a zero of f in G. This means that " is not an isolated zero, so by Theorem 23.3, there must be an open disk D about ", in which f
z is identically zero (Figure 23.4). We want to prove that this forces f
z = 0 for all z in G. To do this, let w be any point of G. We will show that f
w = 0. Since G is connected, there is a path $ in G from " to w. Choose a number r such that every point of $ is at a distance at least 2r from the boundary of G, and also choose r to be less than the radius of D. Now walk along $ from " to w, along the way selecting points at distance less than r from each other. This yields points " = 0  1      n = w on $, as in Figure 23.5. Form an open disk Dj of radius r about each j . (By choice of r, none of these disks reaches the boundary of G). Each j−1 is in Dj−1 , Dj and Dj+1 for j = 1     n − 1. Further, 0 = " is in D0 and D1 , and w is in Dn−1 and Dn . Since 1 is in D0 and D1 , there is a sequence of points in both D0 and D1 converging to 1 (Figure 23.6). But f
z is identically zero on D0 , so f
z vanishes on this sequence. Since this sequence is also in D1 , f
z = 0 for all z in D1 . Now 2 is in D1 and D2 . Choose a sequence of points common to both of these disks, and converging to 2 . Since f
z is identically zero on D1 , then f
z = 0 at each point of this sequence. But since this sequence is also in D2 , then f
z is identically zero on D2 . Continuing in this way, walk along $ from " to w. We find that f
z is identically zero on each of the disks along the way. Finally, f
z is zero on Dn . But w is in Dn , so f
w = 0, and therefore f
z = 0 for all z in G. Proof

This theorem leads immediately to the conclusion that two differentiable functions that agree on a convergent sequence in a domain must be the same function. This is called the identity theorem. COROLLARY 23.1

Identity Theorem

Let f and g be differentiable on a domain G. Suppose f
z and g
z agree on a convergent sequence of distinct points of G. Then f
z = g
z for all z in G. Proof

Apply Theorem 23.4 to the differentiable function f − g.

To get some idea of the power of this result, consider the problem of defining the complex sine function sin
z so that it agrees with the real sine function when z is real. How many ways can this be done?

1016

CHAPTER 23

Series Representations of Functions

Put another way, is it possible to invent two distinct differentiable complex functions, f and g, defined for all z, such that, when z is real, f
x = g
x = sin
x? If this could be done, then we would have f
z = g
z on a convergent sequence of complex numbers (chosen along the real line) in a domain (the entire plane), so necessarily f = g. There can be only one extension of a differentiable function from the real to the complex domain. This is why, when we extend a real function (such as the exponential or trigonometric functions) to the complex plane, we can be assured that this extension is unique.

23.1.2 The Maximum Modulus Theorem Suppose f  S → , and S is a compact set. We know from Theorem 21.3 that f
z assumes a maximum value on S. This means that, for at least one " in S, f
z ≤ f
" for all " in S. But this does not give any information about where in S the point " might be. We will now prove that any such " must lie on the boundary of S, if f is a differentiable function. This is called the maximum modulus theorem. The name of the theorem derives from the fact that the real-valued function f
z is called the modulus of f
z, and we are concerned with the maximum that the modulus of f
z has as z varies over a set S. We will first show that a differentiable function that is not constant on an open disk cannot have its maximum modulus at the center of the disk. LEMMA 23.1

Let f be differentiable and not constant on an open disk D centered at z0 . Then, for some z in this disk, f
z > f
z0   Suppose instead that f
z ≤ f
z0  for all z in D. We will derive a contradiction. Let !
t = z0 + reit for 0 ≤ t ≤ 2. Suppose r is small enough that this circle is contained in D. By Cauchy’s integral theorem, 1 - f
z 1
2 f
z0  = dz = f
z0 + reit dt 2i ! z − z0 2 0 Proof

Then f
z0  ≤

 1
2  f
z0 + reit  dt 2 0

But z0 + reit is in D for 0 ≤ t ≤ 2, so f
z0 + reit  ≤ f
z0 . Then  1
2  1
2 f
z0 + reit  dt ≤ f
z0  dt = f
z0   2 0 2 0 The last two inequalities imply that  1
2  f
z0 + reit  dt = f
z0   2 0 But then

  1
2  f
z0  − f
z0 + reit  dt = 0 2 0

23.1 Power Series Representations

1017

This integrand is continuous and nonnegative for 0 ≤ t ≤ 2. If it were positive for any t, then there would be some subinterval of 0 2 on which the integrand would be positive, and then this integral would be positive, a contradiction. Therefore the integrand must be identically zero:   f
z0 + reit  = f
z0  for 0 ≤ t ≤ 2 This says that f
z has the constant value f
z0  on every circle about z0 and contained in D. But every point in D is on some circle about z0 and contained in D. Therefore f
z = f
z0  = constant for all z in D. Then by Theorem 21.7, f
z = constant on D. This contradiction proves the lemma. We can now derive the maximum modulus theorem. THEOREM 23.6

Maximum Modulus Theorem

Let S be a compact, connected set of complex numbers. Let f be continuous on S and differentiable at each interior point of S. Then f
z achieves its maximum value at a boundary point of S. Further, if f is not a constant function, then f
z does not achieve its maximum at an interior point of S. Because S is compact and f is continuous, we know from Theorem 21.3 that f
z achieves a maximum value at some (perhaps many) points of S. Let " be such a point. If " is an interior point, then there is an open disk D about " that contains only points of S. But then f
z achieves its maximum on this disk at its center. Now there are two cases. Proof

Case 1 f
z is constant on this disk. By the identity theorem, f
z is constant on S. In this event f
z is constant on S. Case 2 f
z is not constant on this disk. Then f
z ≤ f
" for z in this disk, contradicting Lemma 23.1. In this case f
z cannot achieve a maximum in the interior of S, and so must achieve its maximum at a boundary point.

EXAMPLE 23.8

Let f
z = sin
z. We will determine the maximum value of f
z on the square 0 ≤ x ≤ , 0 ≤ y ≤ . First, it is convenient to work with f
z2 , since this will have its maximum at the same value of z that f
z does. Now f
z = sin
z = sin
x cosh
y + i cos
x sinh
y so f
z2 = sin2
x cosh2
y + cos2
x sinh2
y By the maximum modulus theorem, f
z2 must have achieve its maximum value (for this square) on one of the sides of the square. Look at each side in turn. On the bottom side, y = 0 and 0 ≤ x ≤ , so f
z2 = sin2
x achieves a maximum value of 1 On the right side, x =  and 0 ≤ y ≤ , so f
z2 = sinh2
y achieves a maximum value of sinh2
. This is because cos2
 = 1 and sinh
y is a strictly increasing function on 0 .

CHAPTER 23

1018

Series Representations of Functions

On the top side of the square, y =  and 0 ≤ x ≤ . Now f
z2 = sin2
x cosh2
 + cos2
x sinh2
. We need to know where this achieves its maximum value for 0 ≤ x ≤ . This is a problem in single-variable calculus. Let g
x = sin2
x cosh2
 + cos2
x sinh2
 Then g

x = 2 sin
x cos
x cosh2
 − 2 cos
x sin
x sinh2
 = sin
2x cosh2
 − sinh2
 = sin
2x This derivative is zero on
0  at x = /2, so this is the critical point of g. Further,  = cosh2
 g 2 At the ends of the interval, we have g
0 = g
 = sinh2
 < cosh2
 Therefore, on the top side of the square, f
z2 achieves the maximum value of cosh2
. Finally, on the left side of the square, x = 0 and 0 ≤ y ≤ , so f
z2 = sinh2
y, with maximum sinh2
 on 0 ≤ y ≤ . We conclude that, on this square, f
z2 has its maximum value of cosh2
, which is the maximum value of f
z2 on the boundary of the square. Therefore f
z has a maximum value of cosh
 on this square.

SECTION 23.1

PROBLEMS

In each of Problems 1 through 12, find the Taylor series of the function about the point. Also determine the radius of convergence and open disk of convergence of the series. 1. cos
2z z = 0 2. e−z  z = −3i 3.

1  4i 1−z

4. sin
z2  0 5.

1 0
1 − z2

1  1 − 8i 6. 2+z 7. z2 − 3z + i 2 − i 8. 1 +

1 i 2 + z2

9.
z − 92  1 + i 10. e − i sin
z 0 z

11. sin
z + i −i

12.

3  −5 z − 4i

13. Suppose f is differentiable in an open disk about zero, and satisfies f


z = 2f
z + 1. Suppose f
0 = 1 and f

0 = i. Find the Maclaurin expansion of f
z. 14. Find the first three terms of the Maclaurin expansion of sin2
z in three ways as follows: (a) First, compute the Taylor coefficients at 0. (b) Find the first three terms of the product of the Maclaurin series for sin
z with itself. (c) Write sin2
z in terms of the exponential function and use the Maclaurin expansion of this function. 15. Show that 1
2 2z cos
 1 2n z = e d 2 2 0 n=0
n! 

Hint: First show that  n 2 1 z zn = ezw dw n! 2i $ n!wn+1

23.2 The Laurent Expansion

1019

for n = 0 1 2    and $ the unit circle about the origin.

17. Find the maximum value of ez  on the square 0 ≤ x ≤ 1 0 ≤ y ≤ .

16. Find the maximum value of cos
z on the square 0 ≤ x ≤  0 ≤ y ≤ .

18. Find the maximum value of sin
z on the rectangle 0 ≤ x ≤ 2 0 ≤ y ≤ 1.

23.2

The Laurent Expansion If f is differentiable in some disk about z0 , then f
z has a Taylor series representation about z0 . If a function is not differentiable at z0 , it may have a different kind of series expansion about z0 , a Laurent expansion. This will have profound implications in analyzing properties of functions and in such applications as evaluating real and complex integrals. First we need some terminology. The open set of points between two concentric circles is called an annulus. Typically an annulus is described by inequalities r < z − z0  < R in which r is the radius of the inner circle and R the radius of the outer circle (Figure 23.7). We allow r = 0 in this inequality, in which case the annulus 0 < z − z0  < R is a punctured disk (open disk with the center removed). We also allow R = . The annulus r < z − z0  < consists of all points outside the inner circle of radius r. An annulus 0 < z − z0  < consists of all complex z except z0 . We can now state the main result on Laurent series.

THEOREM 23.7

Let 0 ≤ r < R ≤ . Suppose f is differentiable in the annulus r < z − z0  < R. Then, for each z in this annulus, f
z =



cn
z − z0 n 

n=−

where, for each integer n, cn =

f
z 1 dz 2i $
z − z0 n+1

and $ is any closed path about z0 lying entirely in the annulus. A typical such $ is shown in Figure 23.8. The series in the theorem, which may include both positive and negative powers of z − z0 , is the Laurent expansion, or Laurent series, for f
z about z0 in the given annulus. This expansion has the appearance ···+

c−2 c + −1 + c0 + c1
z − z0  + c2
z − z0 2 + · · ·  2
z − z0  z − z0

The function need not be differentiable, or even defined, at z0 or at other points within the inner circle of the annulus. The numbers cn are the Laurent coefficients of f about z0 . A Laurent series is a decomposition of f
z into a sum f
z =

−1  n=−

cn
z − z0 n +

 n=0

cn
z − z0 2 =



 c−n + cn
z − z0 n = h
z + g
z n n=1
z − z0  n=0

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CHAPTER 23

Series Representations of Functions

$ z0 r

R

z0

FIGURE 23.7

FIGURE 23.8

Circles z − z0  = r and z − z0  = R bounding the open annulus r < z − z0  < R.

Closed path $ enclosing z0 and lying in the annulus r < z − z0  < R.

The part containing only nonnegative powers of z − z0 defines a function g
z that is differentiable on z − z0  < R (because this part is a Taylor expansion). The part containing only negative powers of z − z0 defines a function h
z that is not defined at z0 . This part determines the behavior of f
z about a point z0 where f is not differentiable. A proof of the theorem is given on the website. As with Taylor series, we rarely compute the coefficients in a Laurent expansion using this integral formula (quite the contrary, we will use one of these coefficients to evaluate integrals). Instead, we use known series and algebraic or analytic manipulations. This requires that we be assured that the Laurent expansion of a function in an annulus about a point is unique, and does not change with the method of derivation. THEOREM 23.8

Let f be differentiable in an annulus r < z − z0  < R. Suppose, for z in this annulus, f
z =



bn
z − z0 n 

n=−

Then the bn
s are the Laurent coefficients cn of f , and this series is the Laurent expansion of f
z in this annulus. Choose ! as a circle about z0 in the annulus. Let k be any integer. Using Theorem 22.6, we get    f
w 1 n 2ick = dw = bn
z − z0  dw k+1
w − z0 k+1 !
w − z0  ! n=−

Proof

=

 n=−

bn

!

1 dw
w − z0 k−n+1

Now, on !, w = z0 + reit for 0 ≤ t ≤ 2, with r the radius of !. Then
2 1 1 dw = ireit dt k−n+1 r k−n+1
eit k−n+1 !
w − z0  0  i
2 i
n−kt 0 if k = n = k−n e dt =  r 2i if k = n 0

(23.2)

23.2 The Laurent Expansion

1021

Thus in equation (23.2), all terms in the last vanish series except the term with n = k, and the equation reduces to 2ick = 2ibk  Hence for each integer k, bk = ck . Here are some examples of Laurent expansions.

EXAMPLE 23.9

e1/z is differentiable in the annulus 0 < z < , the plane with the origin removed. Since ez =

 1 n z  n=0 n!

then, in this annulus, e1/z =

 1 n=0 n!

 n 1 1 1 1 1 1 1 1 = 1+ + 2 + 3 + +···  z z 2z 6z 24 z4

This is the Laurent expansion of e1/z about 0, and it converges for all nonzero z. This expansion contains a constant term and infinitely many negative integer powers of z, but no positive powers.

EXAMPLE 23.10

We will find the Laurent expansion of cos
z/z5 about zero. For all z, cos
z =


−1n n=0


2n!

z2n 

For z = 0, cos
z  1 1
−1n 2n−5 1 1 1 1 1 = = 5− 3+ z − z+ z3 − · · ·  5 z
2n! z 2 z 24 z 720 40 320 n=0

This is the Laurent expansion of cos
z/z5 about 0. This expansion contains exactly three terms containing negative powers of z, and the remaining terms contain only positive powers. We can think of cos
z/z5 = h
z + g
z, where g
z = −

1 1 z+ z3 − · · · 720 40 320

is a differentiable function (it is a power series about the origin), and h
z =

1 1 1 1 1  − 3+ 5 z 2z 24 z

It is h
z that determines the behavior of cos
z/z5 near the origin.

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CHAPTER 23

Series Representations of Functions

EXAMPLE 23.11

Find the Laurent expansion of 1
z + 1
z − 3i about −1. Use partial fractions to write 1 −1 + 3i 1 1 − 3i 1 = + 
z + 1
z − 3i 10 z + 1 10 z − 3i 1/
z + 1 is already expanded around −1, so concentrate on the last term:   z+1 n 1 1 1 −1  1 = = = z+1 z − 3i −1 − 3i +
z + 1 −1 − 3i 1 − 1+3i 1 + 3i n=0 1 + 3i =−



1
z + 1n 
1 + 3in+1 n=0

This expansion is valid for 
z + 1/
1 + 3i < 1, or z + 1 < 1/
z + 1
z − 3i about −1 is

√ 10. The Laurent expansion of

1 −1 + 3i 1 1 − 3i  1 = −
z + 1n 
z + 1
z − 3i 10 z + 1 10 n=0
1 + 3in+1 √ and this representation is valid in the annulus 0 < z + 1 < 10. √ Notice that 10 is the distance from −1, the center of the Laurent expansion, to the other point, 3i, at which the function is not differentiable.

In the next chapter we will use the Laurent expansion to develop the powerful residue theorem, which has many applications, including the evalution of real and complex integrals.

SECTION 23.2

PROBLEMS

In each of Problems 1 through 10, write the Laurent expansion of the function in an annulus 0 < z − z0  < R about the point. 1. 2. 3. 4. 5.

2z i 1 + z2 sin
z 0 z2 1 − cos
2z 0 z2 z2 cos
i/z 0 z2 1 1−z

6.

z2 + 1  1/2 2z − 1

7.

ez 0 z2

8.

sin
4z 0 z

2

z+i i z−i   1 10. sinh 3  0 z 9.

CHAPTER

24

SINGULARITIES THE RESIDUE THEOREM SOME APPLICATIONS OF THE RESIDUE THEOREM THE ARGUMENT PRINCIPLE AN INVERSION FORMULA FOR THE LAPLACE TRANSFORM EVALUATION OF

Singularities and the Residue Theorem

As a prelude to the residue theorem, we will use the Laurent expansion to classify points at which a function is not differentiable.

24.1

Singularities DEFINITION 24.1

Isolated Singularity

A complex function f has an isolated singularity at z0 if f is differentiable in an annulus 0 < z − z0  < R, but not at z0 itself.

For example, 1/z has an isolated singularity at 0, and sin
z/
z −  has an isolated singularity at . We will now identify singularities as being of different types, depending on the terms appearing in the Laurent expansion of the function about the singularity.

DEFINITION 24.2

Classification of Singularities

Let f have an isolated singularity at z0 . Let the Laurent expansion of f
z in an annulus 0 < z − z0  < R be f
z =



cn
z − z0 n 

n=−

1023

1024

CHAPTER 24

Singularities and the Residue Theorem

Then: 1. z0 is a removable singularity of f if cn = 0 for n = −1 −2     2. z0 is a pole of order m (m a positive integer) if c−m = 0 and c−m−1 = c−m−2 = · · · = 0 3. z0 is an essential singularity of f if c−n = 0 for infinitely many positive integers n.

These three types cover all the possibilities for an isolated singularity. In the case of a removable singularity, the Laurent expansion has no negative powers of z − z0 , and is therefore f
z =



cn
z − z0 n 

n=0

a power series about z0 . In this case we can assign f
z0  the value c0 to obtain a function that is differentiable in the open disk z − z0  < r.

EXAMPLE 24.1

Let f
z =

1 − cos
z z

for 0 < z < . Since cos
z = 1 −

z2 z4 z6 + − +··· 2! 4! 6!

for all z, then f
z =

1 − cos
z z z3 z5 
−1n+1 2n−1 = − + = z z 2! 4! 6! n=1
2n!

for z = 0. The series on the right is actually a power series, having the value 0 at z = 0. We can therefore define a new function 
1 − cos
z/z for z = 0 g
z = 0 for z = 0 which agrees with f
z for z = 0, but is defined at 0 in such a way as to be differentiable there, because g
z has a power series expansion about 0. Because it is possible to extend f to a function g that is differentiable at 0, we call 0 a removable singularity of f . Thus, a removable singularity is one that can be “removed” by appropriately assigning the function a value at the point.

24.1 Singularities

1025

EXAMPLE 24.2

f
z = sin
z/
z −  has a removable singularity at . To see this, first write the Laurent expansion of f
z in 0 < z −  < . An easy way to do this is to begin with sin
z −  = sin
z cos
 − cos
z sin
 = − sin
z so sin
z = − sin
z −  =


−1n+1
z − 2n+1  n=0
2n + 1!

Then, for z = ,
−1n+1 sin
z  1 1 =
z − 2n = −1 +
z − 2 −
z − 4 + · · ·  z −  n=0
2n + 1! 6 120

Although f
 is not defined, the series on the right is defined for z = , and is equal to −1 there. We therefore extend f to a differentiable function g defined over the entire plane by assigning the new function the value −1 when z = :  f
z for z =  g
z =  −1 for z =  This extension “removes” the singularity of f at , since f
z = g
z for z = , and g
 = −1. For f to have a pole at z0 , the Laurent expansion of f about z0 must have terms with negative powers of z − z0 , but only finitely many such terms. If the pole has order m, then this Laurent expansion has the form f
z =

 c−m c−m+1 c−1 + + · · · + + c
z − z0 n 
z − z0 m
z − z0 m−1 z − z0 n=0 n

with c−m = 0. This expansion is valid in some annulus 0 < z − z0  < R.

EXAMPLE 24.3

Let f
z = 1/
z + i. This function is its own Laurent expansion about −i, and c−1 = 1, while all other coefficients are zero. Thus −i is a pole of order 1 of f . This singularity is not removable. There is no way to assign a value to f
−i so that the extended function is differentiable at −i.

EXAMPLE 24.4

g
z = 1/
z + i3 , then g has a pole of order 3 at −i. Here the function is its own Laurent expansion about −i, and the coefficient of 1/
z + i3 is nonzero, while all other coefficients are zero.

1026

CHAPTER 24

Singularities and the Residue Theorem

DEFINITION 24.3

Simple and Double Poles

A pole of order 1 is called a simple pole. A pole of order 2 is a double pole.

EXAMPLE 24.5

Let f
z =

sin
z  z3

For z = 0, f
z = =

1 
−1n 2n+1 
−1n 2n−2 z z = z3 n=0
2n + 1! n=0
2n + 1!

1 2 1 1 1 z − z4 + · · ·  − + 2 z 6 120 5 040

Therefore f has a double pole at 0.

EXAMPLE 24.6

e1/z is defined for all nonzero z, and, for z = 0, e1/z =

 1 1  n! zn n=0

Since this Laurent expansion has infinitely many negative powers of z, 0 is an essential singularity of e1/z . We will discuss several results that are useful in identifying poles of a function.

THEOREM 24.1

Condition for a Pole of Order m

Let f be differentiable in the annulus 0 < z − z0  < R. Then f has a pole of order m at z0 if and only if lim
z − z0 m f
z

z→z0

exists finite and is nonzero. Proof

Expand f
z in a Laurent series in this annulus: f
z =



cn
z − z0 n

for 0 < z − z0  < R

n=−

Suppose f has a pole of order m at z0 . Then c−m = 0 and c−m−1 = c−m−2 = · · · = 0, so the Laurent series is  f
z = cn
z − z0 n  n=−m

24.1 Singularities

1027

Then
z − z0 m f
z =



cn
z − z0 n+m =

n=−m



cn−m
z − z0 n

n=0

= c−m + c−m+1
z − z0  + c−m+2
z − z0 2 + · · ·  Then lim
z − z0 m f
z = c−m = 0.

z→z0

Conversely, suppose limz→z0
z − z0 m f
z = L = 0. We want to show that f has a pole of order m at z0 . Let  > 0. Because of this limit, there is a positive  < R such that 
z − z0 m f
z − L < 

if 0 < z − z0  < 

Then, for such z, 
z − z0 m f
z < L +  In particular, if z − z0  = , then   
z − z0 −n−1 f
z <
L +  z − z0 −n−m−1 =
L +  −n−m−1  The coefficients in the Laurent expansion of f
z about z0 are given by f
z 1 cn = dz 2i $
z − z0 n+1 in which we can choose $ to be circle of radius  about z0 . Then cn  ≤

  1
2max f
z
z − z0 −n−1  < 
L +  −n−m−1 =
L +  −n−m  z on $ 2

Now −n−m can be made as small as we like by choosing  small, if n < −m. We conclude that cn  = 0, hence cn = 0 , if n < −m. Thus the Laurent expansion of f
z about z0 has the form f
z =

 c−m+1 c c−m + + · · · + −1 + cn
z − z0 n  m m−1
z − z0 
z − z0  z − z0 n=0

and therefore f has a pole of order m at z0 , as we wanted to show.

EXAMPLE 24.7

Look again at Example 24.3. Since lim
z + if
z = lim
z + i

z→−i

z→−i

1 = 1 = 0 z+i

f has a simple pole at −i. In Example 24.4, lim
z + i3 g
z = lim
z + i3

z→−i

so g has a pole of order 3 at −i.

z→−i

1 = 1 = 0
z + i3

1028

CHAPTER 24

Singularities and the Residue Theorem

In Example 24.5, lim z2 z→0

sin
z sin
z = 1 = 0 = lim 3 z→0 z z

so sin
z/z3 has a double pole at 0. It is a common error to think that this function has a pole of order 3 at zero, because the denominator has a zero of order 3 there. However, lim z3 z→0

sin
z = lim sin
z = 0 z→0 z3

so by Theorem 24.1, the function cannot have a third-order pole at 0. If f
z is a quotient of functions, it is natural to look for poles at places where the denominator vanishes. Our first result along these lines deals with a quotient in which the denominator vanishes at z0 , but the numerator does not. Recall that g
z has a zero of order k at z0 if g
z0  = · · · = g
k−1
z0  = 0, but g
k
z0  = 0. The order of the zero is the order of the lowest-order derivative that does not vanish at the point.

THEOREM 24.2

Let f
z = h
z/g
z, where h and g are differentiable in some open disk about z0 . Suppose h
z0  = 0, but g has a zero of order m at z0 . Then f has a pole of order m at z0 . We leave a proof of this result to the student.

EXAMPLE 24.8

f
z =

1 + 4z3 sin6
z

has a pole of order 6 at 0, because the numerator does not vanish at 0 , and the denominator has a zero of order 6 at 0. By the same token, f has a pole of order 6 at n for any integer n. Theorem 24.2 does not apply if the numerator also vanishes at z0 . The example f
z = sin
z/z3 is instructive. The numerator has a zero of order 1 at 0, and the denominator a zero of order 3 at 0, and we saw in Example 24.5 that the quotient has a pole of order 2. It would appear that the orders of the zeros of numerator and denominator subtract (or cancel) to give the order of a pole at the point. This is indeed the case.

THEOREM 24.3

Poles of Quotients

Let f
z = h
z/g
z, and suppose h and g are differentiable in some open disk about z0 . Let h have a zero of order k at z0 and g a zero of order m at z0 , with m > k. Then f has a pole of order m − k at z0 . A proof of this is left to the student. By allowing k = 0, this theorem includes the case that the numerator h
z has no zero at z0 .

24.1 Singularities

1029

EXAMPLE 24.9

Consider f
z =


z − 3/24  cos7
z

The numerator has a zero of order 4 at 3/2, and the denominator has a zero of order 7 there, so the quotient f has a pole of order 3 at 3/2.

EXAMPLE 24.10

Let f
z = tan3
z/z9 . The numerator has a zero of order 3 at 0, and the denominator has a zero of order 9 at 0. Therefore f has a pole of order 6 at 0. There are also some useful results stated in terms of products, rather than quotients. We claim that the order of a pole of a product is the sum of the orders of the poles of the factors at a given point.

THEOREM 24.4

Poles of Products

Let f have a pole of order m at z0 and let g have a pole of order n at z0 . Then fg has a pole of order m + n at z0 .

EXAMPLE 24.11

Let f
z =

1  cos4
z
z − /22

Here f
z is a product, which we write for emphasis as

1 1 f
z =  cos4
z
z − /22 Now 1/ cos4
z has a pole of order 4 at /2, and 1/
z − /22 has a pole of order 2 there, so f has a pole of order 6 at /2. f also has poles of order 4 (not 6) at z =
2n + 1/2 for n any nonzero integer other than −1. We are now prepared to develop the powerful residue theorem.

SECTION 24.1

PROBLEMS

In each of Problems 1 through 12, determine all singularities of the function and classify each singularity as removable, a pole of a certain order, or an essential singularity.

1. cos
z/z2 2.

4 sin
z + 2
z + i2
z − i

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CHAPTER 24

Singularities and the Residue Theorem

3. e1/z
z + 2i

10. tan
z

sin
z z− cos
2z 5.
z − 12
1 + z2  z 6.
z + 12

11. 1/ cos
z

4.

12. e1/z
z+1 13. Let f be differentiable at z0 and let g have a pole of order m at z0 . Let f
z0  = 0. Prove that fg has a pole of order m at z0 . 14. Let h and g be differentiable at z0 , g
z0  = 0, and h have a zero of order 2 at z0 . Prove that g
z/h
z has a pole of order 2 at z0 

z−i z2 + 1 sin
z 8. sinh
z z 9. 4 z −1

7.

24.2

15. Suppose h and g are differentiable at z0 and g
z0  = 0, while h has a zero of order 3 at z0 . Prove that g
z/h
z has a pole of order 3 at z0 .

The Residue Theorem To see a connection between Laurent series and the integral of a function, suppose f has a Laurent expansion f
z =



cn
z − z0 n

n=−

in some annulus 0 < z − z0  < R. Let $ be a closed path in this annulus and enclosing z0 . According to Theorem 23.6, the Laurent coefficients are given by an integral formula. In particular, the coefficient of 1/
z − z0  is c−1 = Therefore

$

1 f
zdz 2i $

f
zdz = 2ic−1 

(24.1)

Knowing this one coefficient in the Laurent expansion yields the value of this integral. This fact gives this coefficient a special importance, so we will give it a name.

DEFINITION 24.4

Residue

 n Let f have an isolated singularity at z0 , and Laurent expansion f
z = n=− cn
z − z0  in some annulus 0 < z − z0  < R. Then the coefficient c−1 is called the residue of f at z0 , and is denoted Re s
f z0 .

We will now extend the idea behind equation (24.1) to include the case that $ may enclose any finite number of points at which f is not differentiable.

24.2 The Residue Theorem

!2 z2

!1 z1 !3

!n

$

z3

zn

1031

FIGURE 24.1

THEOREM 24.5

Residue Theorem

Let $ be a closed path and let f be differentiable on $ and all points enclosed by $, except for z1      zn , which are all the isolated singularities of f enclosed by $. Then $

f
zdz = 2i

n 

Re s
f zj 

j=1

In words, the value of this integral is 2i times the sum of the residues of f at the singularities of f enclosed by $. Enclose each singularity zj with a closed path !j (Figure 24.1) so that each !j is in the interior of $, encloses exactly one singularity, and does not intersect any other !k . By the extended deformation theorem, n n   f
zdz = f
zdz = 2i Re s
f zj 

Proof

$

j=1 !j

j=1

The residue theorem is only as effective as our efficiency in evaluating residues of a function at singularities. If we had to actually write the Laurent expansion of f about each singularity to pick off the coefficient of the 1/
z − zj  term, the theorem would be difficult to apply in many instances. What adds to its importance is that, at least for poles, there are efficient ways of calculating residues. We will now develop some of these. THEOREM 24.6

Residue at a Simple Pole

If f has a simple pole at z0 , then Re s
f z0  = lim
z − z0 f
z z→z0

Proof

If f has a simple pole at z0 , then its Laurent expansion about z0 is f
z =

 c−1 + cn
z − z0 n z − z0 n=0

in some annulus 0 < z − z0  < R. Then
z − z0 f
z = c−1 +

 n=0

cn
z − z0 n+1 

1032

CHAPTER 24

Singularities and the Residue Theorem

so lim
z − z0 f
z = c−1 = Re s
f z0 

z→z0

EXAMPLE 24.12

f
z = sin
z/z2 has a simple pole at 0, and sin
z sin
z = 1 = lim z→0 z2 z If $ is any closed path in the plane enclosing the origin, then by the residue theorem, - sin
z dz = 2i Re s
f 0 = 2i z2 $ Re s
f 0 = lim z z→0

EXAMPLE 24.13

Let f
z =

z − 6i 
z − 22
z + 4i

Then f has a simple pole at −4i and a double pole at 2. Theorem 24.6 will not help us with the residue of f at 2, but at the simple pole, Re s
f −4i = lim
z + 4i z→−4i

z − 6i z − 6i −4i − 6i = lim =
z − 22
z + 4i z→−4i
z − 22
−4i − 22

2 3 = − + i 5 10 Before looking at residues at poles of order greater than 1, the following version of Theorem 24.6 is sometimes handy. COROLLARY 24.1

Let f
z = h
z/g
z, where h is continuous at z0 and h
z0  = 0. Suppose g is differentiable at z0 and has a simple zero there. Then f has a simple pole at z0 and Re s
f z0  = Proof

h
z0   g

z0 

f has a simple pole at z0 by Theorem 24.2. By Theorem 24.6, Re s
f z0  = lim
z − z0  z→z0

h
z  h
z h
z = lim =
0  g
z z→z0

g
z − g
z0 /
z − z0  g
z0 

EXAMPLE 24.14

Let f
z =

4iz − 1  sin
z

24.2 The Residue Theorem

1033

Then f has a simple pole at , and, by Corollary 24.1, Re s
f  =

4i − 1 = 1 − 4i cos


In fact, f has a simple pole at n for any integer n, and Re s
f n =

4in − 1 =
−1n
−1 + 4ni cos
n

EXAMPLE 24.15

Evaluate

- 4iz − 1 dz $ sin
z

with $ the closed path of Figure 24.2. y




π

0

π

2π

x

FIGURE 24.2 $ encloses only the singularities −, 0,  and 2 of 4iz−1 . sin
z

$ encloses the poles 0, , 2 and −, but no other singularities of f . By the residue theorem and Example 24.14, - 4iz − 1 dz = 2i Re s
f 0 + Re s
f  + Re s
f 2 + Re s
f − $ sin
z = 2i −1 +
1 − 4i +
−1 + 8i +
1 + 4i = −16 2  Here is a formula for the residue of a function at a pole of order greater than 1. THEOREM 24.7

Residue at a Pole of Order m

Let f have a pole of order m at z0 . Then Re s
f z0  =

dm−1 1 lim m−1
z − z0 m f
z
m − 1! z→z0 dz

1034

CHAPTER 24

Singularities and the Residue Theorem

If m = 1 (simple pole), then
m−1! = 0! = 1 by definition, and the
m−1–order derivative is defined to be just the function itself. With these conventions, the conclusion of the theorem reduces to the result for residues at simple poles when m = 1. Proof

In some annulus about z0 , f
z =

 c−m c−m+1 c + + · · · + −1 + cn
z − z0 n  m m−1
z − z0 
z − z0  z − z0 n=0

It is c−1 we want. Write
z − z0 m f
z = c−m + c−m+1
z − z0  + · · · + c−1
z − z0 m−1 +



cn
z − z0 n+m 

n=0

The right side of this equation is a power series about z0 , and can be differentiated any number of times within its open disk of convergence. Compute dm−1


z − z0 m f
z dzm−1  =
m − 1!c−1 +
n + m
n + m − 1 · · ·
n + 1
z − z0 n+1  n=0

In the limit as z → z0 , this equation yields lim

z→z0

dm−1


z − z0 m f
z =
m − 1!c−1 =
m − 1!Re s
f z0  dzm−1

EXAMPLE 24.16

Let f
z =

cos
z 
z + i3

Then f has a pole of order 3 at −i. By Theorem 24.7,   1 d2 cos
z lim 2
z + i3 Re s
f −i = 2! z→−i dz
z + i3 =

1 1 d2 1 lim 2 cos
z = − cos
−i = − cos
i 2 z→−i dz 2 2

Here are some examples of the residue theorem in evaluating complex integrals.

EXAMPLE 24.17

Let f
z = We want to evaluate of f .

/ $

2iz − cos
z  z3 + z

f
zdz, with $ a closed path that does not pass through any singularity

24.2 The Residue Theorem

1035

The singularities of f are simple poles at 0, i and −i. First compute the residue of f at each of these points. Here it is convenient to use Corollary 24.1: Re s
f 0 =

− cos
0 = −1 1

Re s
f i =

1 2i2 − cos
i −2 − cos
i = = 1 + cos
i 3
i2 + 1 −2 2

and Re s
f −i =

2i
−i − cos
−i 1 = −1 + cos
i 3
−i2 + 1 2

Now consider cases. 1. If $ does not enclose any of the singularities, then

/ $

f
zdz = 0 by Cauchy’s theorem.

2. If $ encloses 0 but not i or −i, then f
zdz = 2i Re s
f 0 = −2i $

3. If $ encloses i but not 0 or −i, then   1 f
zdz = 2i 1 + cos
i  2 $ 4. If $ encloses −i but not 0 or i, then   1 f
zdz = 2i −1 + cos
i  2 $ 5. If $ encloses 0 and i but not −i, then   1 f
zdz = 2i −1 + 1 + cos
i = i cos
i 2 $ 6. If $ encloses 0 and −i but not i, then     1 1 f
zdz = 2i −1 − 1 + cos
i = 2i −2 + cos
i  2 2 $ 7. If $ encloses i and −i but not 0, then   1 1 f
zdz = 2i 1 + cos
i − 1 + cos
i = 2i cos
i 2 2 $ 8. If $ encloses all three singularities, then   1 1 f
zdz = 2i −1 + 1 + cos
i − 1 + cos
i = 2i
−1 + cos
i  2 2 $

EXAMPLE 24.18

Let f
z = We want to evaluate

/ $

sin
z  z2
z2 + 4

f
zdz, where $ is a closed path enclosing 0 and 2i but not −2i.

CHAPTER 24

1036

Singularities and the Residue Theorem

By Theorem 24.3, f has a simple pole at 0, not a double pole, because sin
z has a simple zero at 0. f also has simple poles at 2i and −2i. Only the poles at 0 and 2i are of interest in using the residue theorem, because $ does not enclose −2i. Compute 1 sin
z 1 =  2 z→0 z z +4 4

Re s
f 0 = lim z f
z = lim z→0

and sin
z

Re s
f 2i = lim
z − 2if
z = lim

z→2i z2
z + 2i

z→2i

Then

$

=

sin
2i i = sin
2i
−4
4i 16

  1 i sin
z dz = 2i + sin
2i  z2
z2 + 4 4 16

EXAMPLE 24.19

We will evaluate

-

e1/z dz $

for $ any closed path not passing through the origin. / There are two cases. If $ does not enclose the origin, then $ e1/z dz = 0 by Cauchy’s theorem. If $ does enclose the origin, then use the residue theorem. We need Re s
e1/z  0. As we found in Example 24.6, 0 is an essential singularity of e1/z . There is no simple general formula for the residue of a function at an essential singularity. However, e1/z =

 1 1 n n=0 n! z

is the Laurent expansion of e1/z about 0, and the coefficient of 1/z is 1. Thus Re s
e1/z  0 = 1 and e1/z dz = 2i $

Next we will look at a variety of applications of the residue theorem.

SECTION 24.2

PROBLEMS

In each of Problems 1 through 16, use the residue theorem to evaluate the integral over the given path. -

1+z dz $ is the circle of radius 7
z − 12
z + 2i about −i 2z 2. dz $ is the circle of radius 3 about 1 2 $
z − i

1.

$

2

3. $


ez /zdz $ is the circle of radius 2 about −3i

- cos
z dz $ is the square of side length 3 and sides 2 $ 4+z parallel to the axes, centered at −2i - z+i 5. dz $ is the square of side length 8 and sides 2 $ z +6 parallel to the axes, centered at the origin

4.

24.3 Some Applications of the Residue Theorem -

z−i dz $ is the circle of radius 1 about the origin 2z + 1

6. $

-

z

7.

sinh2
z

$

8.

dz $ is the circle of radius 1 about

1 2

- cos
z i 1 dz $ is the circle of radius about zez 2 8 $ -

iz dz $ is the circle of radius 2 9. 2 + 9
z − i
z $ about −3i 2 10. e2/z dz $ is the square with sides parallel to the $

axes and of length 3, centered at −i 11.

- 8z − 4i + 1 dz $ is the circle of radius 2 about −i z + 4i $

16.

-  z 2 dz $ is any closed path enclosing 1 z−1 $

17. With h and g as in Problem 14 of Section 24.1, show that   2 g

z0 h
4
z0  g


z  g
z Re s  z0 = 2
3 0 −  h
z h
z0  3 h
3
z0 2 18. With h and g as in Problem 15 of Section 24.1, show that   g
z 3 g
z0 h
5
z0  g


z   z0 = 3


0 − Re s h
z h
z0  10
h



z0 2 
4  h
z0  g
z0 h
4
z0  g

z0 h



z0  − +9 24 6
h



z0 3 19. Let g and h be differentiable at z0 . Suppose g
z0  = 0 and let h have a zero of order k at z0 . Prove that g
z/h
z has a pole of order k at z0 , and

-

12.

z2 dz $ is the square of side length 4 and $ z − 1 + 2i sides parallel to the axes, centered at 1 − 2i coth
zdz $ is the circle of radius 2 about i

13. $

14.

-
1 − z2 dz $ is the circle of radius 2 about 2 3 $ z −8 -

15.

e2z dz $ is any closed path enclosing 0 $ z
z − 4i and 4i

24.3

1037

 Re s

k   k! g
z  z0 = h
z h
k
z0 

  Hk   Hk+1   ×  Hk+2       H 2k−1

0 Hk Hk+1   H2k−2

0 0 Hk   H2k−3

··· ··· ···   ···

0 0 0   Hk+1

Gj =

g
j
z0   j!

G0 G1 G2   Gk−1

          

where Hj =

h
j
z0  j!

and

Some Applications of the Residue Theorem 24.3.1

The Argument Principle

The argument principle is an integral formula for the difference between the number of zeros and the number of poles of a function (counting multiplicities) enclosed by a given closed path $. THEOREM 24.8

Argument Principle

Let f be differentiable on a closed path $ and at all points in the set G of points enclosed by $, except at possibly a finite number of poles of f in G. Let Z be the number of zeros of f in G, and P the number of poles of f in G, with each pole and zero counted k times if its multiplicity is k. Then, - f

z dz = Z − P $ f
z Observe first that the only points in G where f
/f might possibly have a singularity are the zeros and poles of f in G.

Proof

1038

CHAPTER 24

Singularities and the Residue Theorem

Now suppose that f has a zero of order k at z0 in G. We will show that f
/f must have a simple pole at z0 , and that Re s
f
/f z0  = k. To see this, first note that, because z0 is a zero of order k, f
z0  = f

z0  = · · · = f
k−1
z0  = 0 while f
k
z0  = 0. Then, in some open disk about z0 , the Taylor expansion of f
z is f
z =



cn
z − z0 n =



cn+k
z − z0 n+k

n=0

n=k

=
z − z0 k



cn+k
z − z0 n =
z − z0 k g
z

n=0

where g is differentiable at z0 (because it has a Taylor expansion there) and g
z0  = ck = 0. Now, in some annulus 0 < z − z0  < R k f

z k
z − z0 k−1 g
z +
z − z0 k g

z g

z = =  + f
z
z − z0 k g
z z − z0 g
z0  Since g

z/g
z is differentiable at z0 , then f

z/f
z has a simple pole at z0 , and Re s
f
/f z0  = k. Next, suppose f has a pole of order m at z1 . In some annulus about z0 , f
z has Laurent expansion f
z =



dn
z − z1 n 

n=−m

with d−m = 0. Then
z − z1 m f
z =



dn
z − z1 n+m =

n=−m



dn−m
z − z1 n = h
z

n=0

with h differentiable at z1 and h
z1  = c−m = 0. Then f
z =
z − z1 −m h
z, so in some annulus about z1 , h

z f

z −m
z − z1 −m−1 h
z +
z − z1 −m h

z −m = =  + −m f
z
z − z1  h
z z − z1 h
z Therefore f
/f has a simple pole at z1 , with Re s
f
/f z1  = −m. Therefore, the sum of the residues of f

z/f
z at singularities of this function in G counts the zeros of f in G, according to multiplicity, and the negative of the number of poles of f in G, again according to multiplicity.

EXAMPLE 24.20

We will evaluate Write

/ $

cot
zdz, with $ the closed path of Figure 24.3. cot
z =

cos
z f

z = sin
z f
z

where f
z = sin
z. Since f has five simple zeros and no poles enclosed by $, the argument principle yields cot
zdz = 2i
5 − 0 = 10i $

24.3 Some Applications of the Residue Theorem

1039

y $


3
2




0



2

3

4

x

FIGURE 24.3

24.3.2

An Inversion Formula for the Laplace Transform

If f is a complex function defined at least for all z on 0 , the Laplace transform of f is
 f
z = e−zt f
tdt 0

for all z such that this integral is defined and finite. If  f = F , then F is the Laplace transform of f , and f is an inverse Laplace transform of F . Sometimes we write f = −1 f, although this requires additional conditions for uniqueness because there are in general many functions having a given F as their Laplace transform. We will give a formula for −1 f in terms of the sum of the residues of ezt F
z at the poles of f . THEOREM 24.9

Inverse Laplace Transform

Let F be differentiable for all z except for a finite number of points z1      zn , which are all poles of F . Suppose for some real , F is differentiable for all z with Re
z > . Suppose also that there are numbers M and R such that zF
z ≤ M

for z > R

For t ≥ 0, let f
t =

n 

Re s
ezt F
z zj 

j=1

Then  f
z = F
z

for Re
z > 

The condition that F is differentiable for Re
z >  means that F

z for all z to the right of the vertical line x = . It is also assumed that zF
z is a bounded function for z outside some sufficiently large circle about the origin. For example, this condition is satisfied by any rational function (quotient of polynomials) in which the degree of the denominator exceeds that of the numerator.

EXAMPLE 24.21

Let a > 0. We want an inverse Laplace transform of F
z = 1/
a2 + z2 . This can be found in tables of Laplace transforms. To use the theorem, F has simple poles at ±ai. Compute  zt  e eati Re s  ai = a2 + z 2 2ai

1040

CHAPTER 24

Singularities and the Residue Theorem

and

 e−ati ezt   −ai = Re s a 2 + z2 −2ai 

An inverse Laplace transform of F is given by f
t =

 1 1  ati e − e−ati = sin
at 2ai a

for t ≥ 0.

EXAMPLE 24.22

We want a function whose Laplace transform is F
z =

1 
z2 − 4
z − 12

F has simple poles at ±2 and a double pole at 1. Compute   ezt ezt 1  2 = lim = e2t  Re s z→2
z + 2
z − 12
z2 − 4
z − 12 4   ezt ezt 1 Re s  −2 = lim = − e−2t  2 2 z→−2
z − 2
z − 12
z − 4
z − 1 36 and

 Re s

 zt   e ezt d  1 = lim 2 2 2 z→1 dz
z − 4
z − 1 z −4 = lim ezt z→1

tz2 − 4t − 2z
z2 − 42

1 2 = − tet − et  3 9

An inverse Laplace transform of F is given by 1 2 1 1 f
t = − tet − et + e2t − e−2t  3 9 4 36 for t > 2 (since all poles of F occur on or to the left of the line Re
z = 2). In these sections we can see a theme developing. A variety of problems (zeros of functions, inverse Laplace transforms, others to be discussed) can be approached by integrating an appropriately chosen complex function over an appropriately chosen path. The function and path must be selected so that the integral gives us the quantity we want to calculate, perhaps after some limit process. We can then use the residue theorem to explicitly evaluate the integral. Depending on the problem, choosing the right function and the right path can be a nontrivial task, but at least this method provides an approach.

24.3.3 Evaluation of Real Integrals We will illustrate the use of the residue theorem in evaluating several general classes of real integrals.

24.3 Some Applications of the Residue Theorem Integrals of for example,

 2 0

K
cos
 sin
d

1041

Let K
x y be a quotient of polynomials in x and y,

x3 y − 2xy2 + x − 2y  x4 + x3 If we replace x with cos
 and y with sin
, we obtain a quotient involving sums of products of positive integer powers of cos
 and sin
. We are interested in evaluating integrals of the form
2 K
cos
 sin
d 0

The idea will be to show that this real integral is equal to an integral of a certain complex function over the unit circle. We then use the residue theorem to evaluate this complex integral, obtaining the value of the real integral. To execute this strategy, let ! be the unit circle, oriented counterclockwise as usual. Parametrize ! by !
 = ei for 0 ≤  ≤ 2. On this curve, z = ei and z = e−i = 1/z, so    1 1  i 1 −i = e +e cos
 = z+ 2 2 z and

   1 1 1  i −i e −e = z−  sin
 = 2i 2i z

Further, on !, dz = iei d = izd so d =

1 dz iz

Now we have     
2 1 1 1 1 1 1 K K
cos
 sin
 i iei d dz = z+  z− 2 z 2i z iz ie 0 !
2 = K
cos
 sin
d 0

This converts the real integral we want to evaluate into the integral of a complex function f
z over the unit circle, where      1 1 1 1 1 f
z = K z+  z−  2 z 2i z iz / Use the residue theorem to evaluate ! f
zdz, obtaining


2 0

K
cos
 sin
d = 2i



Re s
f p

(24.2)

p

The sum on the right is over all of the poles p of f
z enclosed by the unit circle. Poles occurring outside the unit circle are not included in the calculation. Finally, equation (24.2) assumes that f
z has no singularities on the unit circle.  2 The procedure for evaluating 0 K
cos
 sin
d, then, is to compute f
z, determine its poles within the unit circle, evaluate the residues there, and apply equation (24.2). This is a

1042

CHAPTER 24

Singularities and the Residue Theorem

very powerful method that often yields closed form evaluations of integrals for which standard techniques of integration from real calculus are inadequate.

EXAMPLE 24.23

We will evaluate


2 0

sin2
 d 2 + cos


The function K in the above discussion is K
x y =

y2  2+x

The first step is to replace x = cos
 with
z + 1/z/2 and y = sin
 with
z − 1/z/2i, and then multiply by 1/iz, to produce the complex function      1 1 1 1 1 f
z = K z+  z− 2 z 2i z iz 

1  2 z− 1 1 i z4 − 2z2 + 1 =  = 2i 1 z 1 2 + 2
z + z  iz 2 z2
z2 + 4z + 1 √ f has√a double pole at 0 and simple poles at zeros of z2 + 4z + 1, which are −2 + 3 and these two simple poles of f , the first is enclosed by ! and the second is not, so −2 − 3. Of √ discard −2 − 3. By equation (24.2),


2 0

 √  sin2
 d = 2i Re s
f 0 + Re s
f −2 + 3  2 + cos


Now  d 2 d i z4 − 2z2 + 1 z f
z = lim z→0 dz z→0 dz 2 z2 + 4z + 1  5  4 i z + 6z − 4z2 − 3z + 2z3 − 2 = lim 2 = −2i 2 z→0
z2 + 4z + 12

Re s
f 0 = lim

and



√ z4 − 2z2 + 1 i Re s
f −2 + 3 = 2 2z
z2 + 4z + 1 + z2
2z + 4 z=−2+√3 √ i 42 − 24 3 = √  2 −12 + 7 3

Then


2 0

 √  √ i 42 − 24 3 sin2
 90 − 52 3 d = 2i −2i + = √ √  2 + cos
 2 −12 + 7 3 12 − 7 3

approximately 168357. In applying this method, if a complex number results, check the calculations, because a real integral must have a real value.

24.3 Some Applications of the Residue Theorem

1043

EXAMPLE 24.24

We will evaluate




 0

1 d  +  cos


where 0 <  < . Since the method we have developed deals with integrals over 0 2, we must first decide how to accommodate an integral over 0 . Write
2
2
 1 1 1 d + d d =  +  cos
  +  cos
 0 0  +  cos
  Let w = 2 −  in the last integral to obtain
0

2 1 1 1 d =
−1dw = dw  +  cos
    +  cos
2 − w 0  +  cos
w Therefore


0



1
2 1 1 d = d  +  cos
 2 0  +  cos


and we can concentrate on the integral over 0 2. First produce the function f
z = f has simple poles at

1 −2i 1    = 2 1 iz z + 2z +  + z+ z  2

 − ± 2 − 2 z=  

Since  > , these numbers are real. Only one of them,  − + 2 − 2 z1 =  is enclosed by !. The other is outside the unit disk and is irrelevant for our purposes. Then
 1
2 1 1 1 d = d = 2i Re s
f z1  2 0  +  cos
 2 0  +  cos
 = i

 −2i =  2 2z1 + 2  − 2

Before continuing with other kinds of real integrals we can evaluate using the residue theorem, we will take a brief excursion and give another, perhaps surprising, proof of the fundamental theorem of algebra. This argument is originally due to N.C. Ankeny, and the version we give appeared in Lion Hunting and Other Mathematical Pursuits, by R.P. Boas (The Mathematical Association of America Dolciani Mathematical Expositions, Volume 15). Let p
z be a nonconstant polynomial with complex coefficients. We want to show that, for some number z, p
z = 0. First, we may assume that p
x is real if x is real. To see why this is true, let p
z = a0 + a1 z + · · · + an zn 

1044

CHAPTER 24

Singularities and the Residue Theorem

where an = 0. Denote p
z = a0 + a1 z + · · · + an zn  Then q
z = p
zp
z is a nonconstant polynomial. Further, if z = x is real, then x = x and q
x = p
xp
x =
a0 + a1 x + · · · + an xn 
a0 + a1 x + · · · + an xn    =
a0 + a1 x + · · · + an xn  a0 + a1 x + · · · + an xn = a0 + a1 x + · · · + an xn 2 is real. We could then use q
z in our argument in place of p
z, since q
z is a polynomial with no zero if p
z has no zero. Thus suppose that p
z = 0 for all z, and p
x is real if x is real. Because p
x is continuous and never zero for real x, p
x must be strictly positive or strictly negative for all real x. But then
2 1 d = 0 p
2 cos
 0 But, by the method we have just discussed, with ! the unit circle, we conclude that
2 1 1 - zn−1 1 1 d = dz = dz = 0 p
2 cos
 i ! r
z ! p
z + 1/z iz 0 where r
z = zn p
z + 1/z      2 n  1 1 1 + a2 z + = zn a0 + a1 z + + · · · + an z + z z z

z2 + 1
z2 + 12
z2 + 1n + a2 = zn a0 + a1 + · · · + a  n z z2 zn From this it is clear that r
z is a polynomial. If r
" = 0 for some " = 0, then we would have p
" + 1/" = 0, so " + 1/" would be a zero of p, a contradiction. Further, r
0 = an = 0 because p has degree n. Therefore r
z = 0 for all z, so zn−1 /r
z is a differentiable function for all z. But then, by Cauchy’s theorem, 1 - zn−1 dz = 0 i ! r
z a contradiction. We conclude that p
z = 0 for some number z, proving the fundamental theorem of algebra. Evaluation of

 −

 p
x/q
xdx

We will now consider real integrals of the form
p
x dx − q
x

in which p and q are polynomials with real coefficients and no common factors, q has no real zeros, and the degree of q exceeds the degree of p by at least 2. These conditions are sufficient to ensure convergence of this improper integral. As with the preceding class of integrals, the strategy is to devise a complex integral that is equal to this real integral, then evaluate the complex integral using the residue theorem. To

24.3 Some Applications of the Residue Theorem

1045

y ! z3
R

z2 zm R

z1

S

R

x

FIGURE 24.4

do this, first observe that q
z has real coefficients, so its zeros occur in complex conjugate pairs. Suppose the zeros of q are z1  z1  z2  z2      zm  zm , with each zj in the upper half plane Im
z > 0 and each zj in the lower half plane Im
z < 0. Let $ be the curve shown in Figure 24.4, consisting of a semicircle ! of radius R and the segment S from −R to R on the real axis, with R large enough that $ encloses all the poles z1      zm of p
z/q
z in the upper half plane. Then
p
z
p
z - p
z m  dz = 2i Re s
p/q zj  = dz + dz (24.3) $ q
z S q
z ! q
z j=1 On S, z = x for −R ≤ x ≤ R, so


R p
x
p
z dz = dx −R q
x S q
z

Next consider the integral over !. Since the degree of q
z exceeds that of p
z by at least 2, degree of z2 p
z ≤ degree of q
z This means that, for sufficiently large R, z2 p
z/q
z is bounded for z ≥ R. That is, for some number M,  2   z p
z     q
z  ≤ M for z ≥ R Then

so

   p
z  M M    q
z  ≤ z2 ≤ R2

for z ≥ R

 
 p
z  M    ! q
z dz ≤ R2
length of !

M M → 0 as R → 
R = R2 R Thus, in the limit as R → in equation (24.3), the first integral on the right has limit 
p
x/q
xdx, and the second integral has limit zero. In the limit as R → , equation − (24.3) yields  
p
x m  p
z dx = 2i Re s  zj  (24.4) q
z − q
x j=1 =

Equation (24.4) provides a general method for evaluating integrals of rational functions over the real line, under the assumptions made above. It is not necessary to repeat the derivation of this equation each time it is used—simply determine the zeros of q
z in the upper half-plane, evaluate the residue of p/q at each such zero (which is a pole of p/q whose order must be determined), and apply equation (24.4).

1046

CHAPTER 24

Singularities and the Residue Theorem

EXAMPLE 24.25

We will evaluate


−

1 dx x6 + 64

Here p
z = 1 and q
z = z6 + 64. The degree of q exceeds that of p by 6, and q has no real zeros. The zeros of z6 + 64 are the sixth roots of −64. To find these, put −64 in polar form: −64 = 64ei
+2n  in which n can be any integer. The six sixth roots of −64 are 2ei
+2n/6 for n = 0 1 2 3 4 5 The three roots in in the upper half-plane are z1 = 2ei/6  z2 = 2ei/2 = 2i

and z3 = 2e5i/6 

We need the residue of 1/
z6 + 64 at each of these simple poles. Corollary 24.1 is convenient to use here:   1 1 −5i/6 1 i/6  2e e Re s 6 =  = i/6 5 z + 64 6
2e  192   1 i 1  2i =  Re s 6 =− z + 64 6
2i5 192 and  Re s

 1 1 −25i/6 1 −i/6 1 5i/6 = =  =  2e e e 6 5i/6 5 z + 64 6
2e  192 192

Then




−

 2i −5i/6 1 dx = − i + e−i/6 e x6 + 64 192        5 5 i = cos − i sin − i + cos − i sin  96 6 6 6 6

Now  cos

  5 + cos =0 6 6

and 

5 sin 6

 = sin

 6

=1

so




−

1 i  dx =
−2i =  x6 + 64 96 48

24.3 Some Applications of the Residue Theorem

1047

  Integrals of −  p
x/q
xcos
cxdx and −  p
x/q
xsin
cxdx Suppose p and q are polynomials with real coefficients and no common factors, that the degree of q exceeds the degree of p by at least 2, and that q has no real zeros. We want to evaluate integrals
p
x
p
x cos
cxdx and sin
cxdx − q
x − q
x in which c is any positive number. Again, we proceed by looking for the integral of a suitably chosen complex function over a suitably chosen closed curve. Consider - p
z eicz dz $ q
z where $ is the closed path of the preceding subsection, enclosing all the zeros z1      zm of q lying in the upper half plane. Here is why this integral is promising. With $ consisting of the semicircle ! and the segment S on the real axis, as before, then
p
z
R p
x - p
z eicz dz = eicz dz + eicx dx $ q
z ! q
z −R q
x
R p
x - p
z
R p
x eicz dz + cos
cxdx + i sin
cxdx = ! q
z −R q
x −R q
x   m  p
z icz = 2i Re s e  zj  q
z j=1  As R → , one can show that ! p
z/q
zeicz dz → 0, leaving  
p
x
p
x m  p
z icz cos
cxdx + i sin
cxdx = 2i Re s e  zj  (24.5) q
z − q
x − q
x j=1  The real partof the right side of equation (24.5) is − p
x/q
x cos
cxdx, and the imaginary part is − p
x/q
x sin
cxdx.

EXAMPLE 24.26

We will evaluate




−

cos
cx dx
x2 + 2 
x2 + 2 

in which c,  and  are positive numbers and  = . The zeros of the denominator in the upper half plane are i and i, and these are simple poles of f
z =

eicz 
z2 + 2 
z2 + 2 

Compute Re s
f i =

eici e−c = 2 2 2i
 −   2i
2 − 2 

and Re s
f i =

e−c  2i
2 − 2 

1048

CHAPTER 24 Then

Singularities and the Residue Theorem


cos
cx sin
cx dx + i dx 2 2 2 2 2 2 2 2 −
x +  
x +   −
x +  
x +    −c

 e−c e e−c e−c  = 2i + − = 2  2i
2 − 2  2i
2 − 2   − 2  




Separating real and imaginary parts, we have  −c 
e  e−c cos
cx dx = 2 − 2 2 2 2  − 2   −
x +  
x +   and






−

sin
cx dx = 0
x2 + 2 
x2 + 2 

The latter is obvious because the integrand is an odd function. Integrals Using Indented Contours Equation (24.4) enables us to evaluate certain improper integrals of quotients of polynomials, assuming that the denominator has no real zeros. We will extend this result to the case that the denominator has simple real zeros. Consider
p
x dx − q
x in which p and q are polynomials with real coefficients and no common factors, and the degree of q exceeds that of p by at least 2. Suppose q has complex zeros z1      zm in the upper half plane, as well as simple real zeros t1      tk . Let $ be the path of Figure 24.5, including a semicircle ! of radius R about the origin, small semicircles !j of radius  centered at each real zero tj , and segments Lj along the real line connecting these semicircles as shown. We call such a path an indented path because of the small semicircles about the real zeros of q
x. Let  be small enough that no two of these semicircles intersects the other semicircles, and no zj is enclosed by any !k with j = k. Also suppose R is large enough that $ encloses each zj . Note that each tj is outside $. By the residue theorem,   - p
z  p
z dz = 2i Re s z q
z j $ q
z j=1 =


p
z k
k+1  
p
x p
z dz + dz + dx ! q
z j=1 !j q
z j=1 Lj q
x

We want to investigate what happens in equation (24.6) when R → and  → 0. y !

!1

!2

L1 t1 L2 t2 FIGURE 24.5

!3 L3

t3 L 4

x

(24.6)

24.3 Some Applications of the Residue Theorem

1049

 When R → , we claim that ! p
z/q
zdz → 0 by an argument like that we have done before. The sum of the residues at zeros of q is unchanged in this limit. As  → 0, the semicircles !j contract to tj , and the segments Lj expand to cover the interval −R R, and then the entire real line as R → . This means that, in equation (24.6) k+1 
j=1 Lj


p
x p
z dz → dx q
z − q
x

 It is not yet clear what happens to each integral !j p
z/q
zdz in this process. We will show that each of these integrals approaches i times the residue of p
z/q
z at the real simple pole tj . To see this, write the Laurent expansion of p
z/q
z about tj :  p
z c c = −1 + cs
z − tj s = −1 + g
z q
z z − tj s=0 z − tj

where g is differentiable at tj . On !j , z = tj + eit , where t varies from  to 0 (for counterclockwise orientation on $). Then


p
z 1 dz = c−1 dz + g
zdz !j q
z !j z − t j !j
0
0 1 it ie dt + g
tj + eit ieit dt = c−1 it  e 
0 = −ic−1 + i g
tj + eit eit dt  = −i Re s Now i

0 




0 p
z  tj + i g
tj + eit eit dt q
z 

g
tj + eit eit dt → 0 as  → 0. Therefore  
p
z p
z dz → −i Re s t  q
z j !j q
z

Therefore, as R → and  → 0 in equation (24.6), we get 2i

 j=1

 Re s

  
k  p
z p
z p
x  zj = −i Re s  tj + dx q
z q
z q
x − j=1

hence




−

    k   p
z p
z p
x dx = i Res  t + 2i Re s z  q
x q
z j q
z j j=1 j=1

(24.7)

In a sense, the simple poles of p
z/q
z on the real line contribute “half residues”, having been enclosed by semicircles instead of circles, while the poles in the upper half plane contribute “full residues” to this sum.

1050

CHAPTER 24

Singularities and the Residue Theorem

EXAMPLE 24.27

Evaluate


−

3x + 2  x
x − 4
x2 + 9

Here f
z =

3z + 2  z
z − 4
z2 + 9

The denominator has simple real zeros at 0 and 4, and simple complex zeros −3i and 3i. Only 3i is in the upper half plane. Compute the residues: 1 2 =−  z→0 −36 18 7 14 Re s
f 4 = lim
z − 4f
z = = z→4 100 50

Re s
f 0 = lim zf
z =

and Re s
f 3i = lim
z − 3if
z = z→3i

Then




−

Integrals

 0

2 + 9i 9i + 2 =  3i
3i − 4
6i 72 − 54i

    1 7 2 + 9i 3x + 2 14 = i − + + 2i = −  2 x
x − 4
x + 9 18 50 72 − 54i 75

xa p
x/q
xdx

Let 0 < a < 1. We will consider integrals of the form


0

xa

p
x dx q
x

in which p and q are polynomials with real coefficients and no common factors, q has no positive zeros, and the degree of q exceeds the degree of p by at least 1. We also assume that either q
0 = 0, or q
z has a simple zero at the origin. Let the nonzero zeros of q be z1      zm . These are all the nonzero zeros of q, not just those in the upper half-plane. Since the coefficients of q are real, this list includes complex conjugate pairs. Choose r small enough and R large enough that the closed path $ shown in Figure 24.6 encloses z1      zm . $ consists of !R (“most” of the circle of radius R about 0), !r (“most” of the circle of radius r about 0), and the line segments L1 and L2 connecting !r and !R . We will eventually let r → 0 and R → , but some work is required first. We must agree on a meaning for za , since this symbol generally denotes a (possibly infinite) set of different numbers. Write z = ei for some  in 0 2, and define za = a eia  As z approaches L1 , f
z =

za p
z xa p
x →  q
z q
x

24.3 Some Applications of the Residue Theorem

1051

y !R z1

z2 !r

L1 x L2

zm z4

z3

FIGURE 24.6

where r < x < R. But as z approaches L2 , the lower side of the positive real axis, f
z →

xa e2ai p
x  q
x

The reason for this is that the argument increases by 2 as z approaches the positive real axis from below, and then za = a ei
+2a = a eia e2ai  By the residue theorem, - za p
z m  dz = 2i Res
f zj  $ q
z j=1 =


!R


za p
z
xa p
x
xa e2ai p
x za p
z dz + dz + dx + dx q
z q
x !r q
z L1 q
x L2

On L1 , x varies from r to R, while on L2 , x varies from R to r to maintain counterclockwise orientation on $. The last equation becomes 2i

m 

Re s
f zj  =


za p
z za p
z dz + dz q
z !r q
z


!R

j=1

+




R r


r xa e2ai p
x xa p
x dx + dx q
x q
x R

By making estimates on the two circular arcs, we can show that the first two integrals in the last equation tend to zero as r → 0 and R → . In this limit, the last equation becomes 2i

m 

Re s
f zj  =

j=1




0


0 xa e2ai p
x xa p
x dx + dx q
x q
x

From this we obtain
0



m 2i  xa p
x dx = Re s
f zj  q
x 1 − e2ai j=1

(24.8)

1052

CHAPTER 24

Singularities and the Residue Theorem

EXAMPLE 24.28

We will use equation (24.8) to evaluate


0

Here p
z = 1, a =

1 3

x1/3 dx x
x2 + 1

and q
z = z
1 + z2 , with simple zeros at 0, i and −i. Compute

Re s
z1/3 /q
z i =

1 i1/3
ei/2 1/3 = − ei/6 = 2i2 −2 2

and Re s
z1/3 q
z −i = Then


0



1
e3i/2 1/3 = − ei/2  −2 2

   2i x1/3 1  i/6 dx = e + ei/2 − 2 2i/3 x
x + 1 1−e 2  √ −i 3 1 √ = + i + i =  2 1+ 1 − 3i 2 2

2

Many other kinds of real integrals can be evaluated using complex integral techniques. Some of these require considerable ingenuity in seeking the right function to integrate over the path to get the result that is wanted. The Cauchy Principal Value Since we have been dealing with improper integrals, we will mention the Cauchy principal value. An integral
I= g
xdx −

is defined to be




0

lim

r→− r

g
xdx + lim




R→ 0

R

g
xdx

if both of these integrals converge. These limits are independent of each other. The Cauchy principal value of I is defined to be  

R CPV g
xdx = lim g
xdx −

R→ −R

This is a special case of  the two independent limits defining I. In the event that − g
xdx converges, certainly the value of I agrees with the Cauchy principal value of the integral. However, it is possible for an integral to  have a finite CPV, but to diverge in the broader sense of the definition of I. This occurs with − xdx, which certainly diverges. However, this integral has Cauchy principal value 0, because for any positive R,
R xdx = 0 −R

In some of the examples we have discussed, we were actually computing Cauchy principal R values, whenever we took a limit of −R g
xdx as R → . In most cases the conditions imposed insured that the improper integral converged in the more general sense as well.

24.3 Some Applications of the Residue Theorem

SECTION 24.3

PROBLEMS

-

z dz, with $ the circle z = 2, first 2 $ 1+z using the residue theorem, and then by the argument principle.

1. Evaluate

/ 2. Evaluate $ tan
zdz, with $ the circle z = , first by the residue theorem, and then by the argument principle. -

z+1 dz, with $ the circle z = 1, 2 $ z + 2z + 4 first by the residue theorem, and then by the argument principle.

3. Evaluate

4. Let p
z =
z − z1 
z − z2  · · ·
z − zn , with z1      zn distinct complex numbers. Let $ be a positively oriented closed path enclosing each of the z
j s. Prove that - p

z dz = 2in $ p
z first by the residue theorem, and then by the argument principle. In each of Problems 5 through 9, find an inverse Laplace transform of the function, using residues.


14.

1
z + 32

7.

1
z − 22
z + 4

8.


15.

In each of Problems 10 through 22, evaluate the integral.


2

1 d 6 + sin


2

1 d 2 − cos


10. 0


11. 0






−


13.



−




2

16.


−


18.

−


19.



−







1 dx x
x + 4
x2 + 16 sin
x dx x2 − 4x + 5 cos2
x dx
x2 + 42



1 dx
x − 4
x5 + 1

−

22.

d

sin
 + cos
 d 2 − cos


0




2 sin


2

20.

21.

x sin
2x dx x4 + 16 2 + sin2


0

17.

1 dx x2 − 2x + 6



x3/4 dx x4 + 1

23. Let  be a positive number. Show that
cos
x dx = e−  2 − x + 1 24. Let  and  be positive numbers. Show that
cos
x  dx = 3
1 + e−  2 2 2 − x + 

1
z2 + 9
z − 22

9.
z + 5−3

12.



−

0

z 5. 2 z +9 6.

1053



−

1 dx x4 + 1 1 dx x6 + 1

25. Let  and  be distinct positive numbers. Show that
2 2 1 d =   0 2 cos2
 + 2 sin2
 26. Let  be a positive number. Show that
/2  1  d =  0  + sin2
 2 
1 +  27. Let  be a positive number. Show that √
 −2 2 e  e−x cos
2xdx = 2 0 2

Hint: Integrate e−z about the rectangular path having corners at ±R and R ± i. Use Cauchy’s theorem to evaluate this integral, set this equal to the sum of the

1054

CHAPTER 24

Singularities and the Residue Theorem

integrals on the sides of the rectangle, and take the limit as R → . Assume the standard result that √
 2  e−x dx = 2 0

y  R

28. Derive Fresnel’s integrals: 

1   cos
x2 dx = sin
x2 dx = 2 2 0 0 2

Hint: Integrate eiz over the closed path bounding the sector 0 ≤ x ≤ R, 0 ≤  ≤ /4, show in Figure 24.7. Use Cauchy’s theorem to evaluate this integral, then evaluate it as the sum of integrals over the boundary segments of the sector. Show that the integral over the circular sector tends to zero as R → , and use the integrals over the straight segments to obtain Fresnel’s integrals.

x

FIGURE 24.7

29. Let  and  be positive numbers. Show that  
x sin
x  −/√2  dx = e sin  √ x4 + 4 22 0 2 30. Let 0 <  < . Show that


 0

 1 d = 2 
 +  cos
2
 − 2 3/2

CHAPTER

25

FUNCTIONS AS MAPPINGS CONFORMAL MAPPINGS CONSTRUCTION OF CONFORMAL MAPPINGS BETWEEN DOMAINS HARMONIC FUNCTIONS AND THE DIRICHLET PROBLEM COMPLEX FUNCTION MODELS

Conformal Mappings

In the calculus of real functions of a single real variable, we may gain some insight into the function’s behavior by sketching its graph. For complex functions we cannot make the same kind of graph, since a complex variable z = x + iy by itself has two variables. However, we can set w = f
z and make two copies of the complex plane, one for z and the other for image points w. As z traces out a path or varies over a set S in the z plane, we can plot image points w = f
z in the w plane, yielding a picture of how the function acts on this path or on points in S. The set of all image points f
z for z in S is denoted f
S. When looked at in this way, a function is called a mapping or transformation. This idea is diagrammed in Figure 25.1. z S

w w  f (z)

f (S)

FIGURE 25.1

Thinking of a function as a mapping can be a powerful tool in solving certain kinds of problems, including the analysis of fluid motion and the solution of partial differential equations, particularly Dirichlet problems. We will develop some ideas about mappings, then turn to applications.

25.1

Functions as Mappings First we need some terminology. Let f be a complex function and D a set of points in the plane on which f
z is defined. Let D∗ also be a set of complex numbers. 1055

1056

CHAPTER 25

Conformal Mappings

DEFINITION 25.1

1. f maps D into D∗ if f
z is in D∗ for every z in D. In this event, we write f  D → D∗ . 2. f maps D onto D∗ if f
z is in D∗ for every z in D and, conversely, if w is in D∗ , then there is some z in D such that w = f
z. In this event, we call f an onto mapping.

Thus f  D → D∗ is onto if every point of D∗ is the image under f of some point in D.

EXAMPLE 25.1

Let f
z = iz for z ≤ 1. Then f acts on points in the closed unit disk D  z ≤ 1. If z is in D, then f
z = iz = z ≤ 1, so the image of any point in this disk is also in this disk. Here f maps D into D. (So D∗ = D in the definition). This mapping is onto. If w is in D, then z = w/i is also in D, and f
z = f
w/i = i
w/i = w Every point in the unit disk is the image of some point in this disk under this mapping. We can envision this mapping geometrically. Since i = ei/2 , if z = rei , then f
z = iz = rei/2 ei = rei
+/2  so f takes z and adds /2 to its argument. This rotates the line from the origin to z by /2 radians counterclockwise. The action of f on z can be envisioned as in Figure 25.2. Since this function is simply a counterclockwise rotation by /2 radians, it is clear why f maps the unit disk onto itself.

z

w i

i

f(z)  iz z 

  2π 1

FIGURE 25.2

w  iz

1

The mapping f
z = iz for z ≤ 1.

Often we have the function f and the set D of complex numbers to which we want to apply this mapping. We then must analyze f
z to determine the image of D under the mapping. In effect, we are finding D∗ so that f is a mapping of D onto D∗ .

EXAMPLE 25.2

Let f
z = z2 for z in the wedge D shown in Figure 25.3. D consists of all complex numbers on or between the nonnegative real axis and the line y = x.

25.1 Functions as Mappings z

w  z2

1057

w D*

D π

π

4

2

FIGURE 25.3 w = z2 maps D

one-to-one onto D∗ .

In polar form z = rei is in D if 0 ≤  ≤ /4. Then f
z = z2 = r 2 e2i , so f has the effect of squaring the magnitude of z, and doubling its argument. If z has argument between 0 and /4, then z2 has argument between 0 and /2. This spreads the wedge D out to cover the entire right quarter-plane, consisting of points on or between the nonnegative real and imaginary axes. If we call this right quarter plane D∗ , then f maps D onto D∗ . Some functions map more than one point to the same image. For example, f
z = sin
z maps all integer multiples of  to zero. If each image point is the unique target of exactly one number, then the mapping is called one-to-one.

DEFINITION 25.2

A mapping f  D → D∗ is one-to-one if distinct points of D map to distinct points in D∗ .

Thus f is one-to-one (or 1 − 1) if z1 = z2 implies that f
z1  = f
z2 . The notions of one-to-one and onto are independent of each other. A mapping may have one of these properties, or both, or neither. The mapping f
z = z2 of Example 25.2 maps the wedge 0 ≤ arg
z ≤ /4 in one-to-one fashion onto the right quarter plane. However, f
z = z2 does not map the entire complex plane to the complex plane in a one-to-one manner, since f
−z = f
z. This function does map the plane onto itself, since, given any complex number w, there is some z such that f
z = z2 = w.

EXAMPLE 25.3

Let h
z = z2 for all z. h maps the entire plane onto itself, but is not one-to-one. If z = x + iy, then h
z = x2 − y2 + 2ixy = u + iv where u = x2 − y2 and v = 2xy. We will use this formulation to determine the image under f of a vertical line x = a. Any point on this line has the form z = a + iy, and maps to h
a + iy = u + iv = a2 − y2 + 2iay Points on the line x = a map to points
u v with u = a2 − y2 and v = 2ay. Write y = v/2a (assuming that a = 0) to obtain u = a2 −

v2  4a2

or v2 = 4a2
a2 − u the equation of a parabola in the u v plane. h maps vertical lines x = a = 0 to parabolas.

1058

CHAPTER 25

Conformal Mappings

If a = 0, the vertical line x = a is the imaginary axis, consisting of points z = iy. Now h
z = −y2 , so h maps the imaginary axis in the x y plane to the nonpositive part of the real axis in the u v plane. Figure 25.4 shows the parabolic image of a line x = a = 0. The larger a is, the more the parabola opens up to the left (intersects the v axis further from the origin). As a is chosen smaller, approaching 0, these parabolas become “flatter”, approaching the nonpositive real axis in the u v plane. A horizontal line y = b consists of points z = x + ib, which map to h
z =
x + ib2 = x2 − b2 + 2ixb = u + iv Now u = x2 − b2 and v = 2xb, so, for b = 0, v2 = 4b2
b2 + u A typical such parabola is shown in Figure 25.5, opening up to the right. These parabolas also upen up more the larger b is. If b = 0, the line y = b is the real axis in the z plane, and this maps to h
x = x2 , giving the nonnegative real axis in the u v plane as x takes on all real values.

z

w  z2

y a

x

w

z v a2

bi

u

Im(z)  b

x

Re z  a FIGURE 25.4 w = z2 maps vertical lines to parabolas opening left.

w

w  z2

y

v


b 2

u

FIGURE 25.5 w = z2 maps horizontal lines to parabolas opening right.

EXAMPLE 25.4

We will look at the exponential function w = E
z = ez as a mapping. Write w = u + iv = ex+iy = ex cos
y + iex sin
y so u = ex cos
y

and v = ex sin
y

As a mapping of the entire plane to itself, E is not onto (no number maps to zero), and E is also not one-to-one (all points z + 2ni have the same image, for n any integer). Consider a vertical line x = a in the x y plane. The image of this line consists of points u + iv with u = ea cos
y

v = ea sin
y

Then u2 + v2 = e2a  so the line x = a maps to the circle of radius ea about the origin in the u v plane. Actually, as the point z = a + iy moves along this vertical line, the image point u + iv makes one complete circuit around the circle every time y varies over an interval of length 2, since cos
y + 2n = cos
y and sin
y + 2n = sin
y. We may therefore think of a vertical line as infinitely many intervals of length 2 strung together, and the exponential function wraps each segment once around the circle u2 + v2 = e2a (Figure 25.6).

25.1 Functions as Mappings z

w

w  ez

y a

z v

ea

x

w 

Re z  a

b

u

w = ez wraps a vertical line around a circle, covering the circle once for every interval of length 2.

Im z  d

d

Im z  c

c

Im z  b

b

x

u

lines to half-rays from the origin.

w

v

b

θd θc

w  ez

Re z  a

v

FIGURE 25.7 w = ez maps horizontal

y

a

w

ea

FIGURE 25.6

z

w  ez

y

1059

u

x w  e a w  e b

Re z  b

FIGURE 25.8 w = ez maps the rectangle shown to a wedge bounded by two half-rays and two circles.

The image of a point z = x + ib on the horizontal line y = b is a point u + iv with u = ex sin
b

y = ex cos
b

As x varies over the real line, ex varies from 0 to over the positive real axis. The point
ex sin
b ex cos
b moves along a half-line from the origin to infinity, making an angle b radians with the positive real axis (Figure 25.7). In polar coordinates this half-line is  = b. Using these results, we can find the image of any rectangle in the x y plane, having sides parallel to the axes. Let the rectangle have sides on the lines x = a, x = b, y = c and y = d (in the x y plane in Figure 25.8). These lines map, respectively, to the circles u2 + v2 = e2a 

u2 + v2 = e2b

and the half-lines =c

and  = d

The wedge in the u v plane in Figure 25.8 is the image of the rectangle under this exponential mapping. Given a mapping f and a domain D, here is a strategy that is often successful in determining f
D. Suppose D has a boundary made up of curves !1      !n . Find the images of these curves, f
!1      f
!n . These form curves in the w− plane, bounding two sets, labeled I and II in Figure 25.9. f
D is one of these two sets. To determine which it is, choose a convenient point " in D and locate f
". This set will be f
D.

EXAMPLE 25.5

We will determine the image, under the mapping w = f
z = sin
z, of the strip S consisting of all z with −/2 < Re
z < /2 and Im
z > 0. S is shown in Figure 25.10.

1060

CHAPTER 25

Conformal Mappings z z !1 !2

w

f(!1)

!n

f(!n)

II

I

D

Re z 
2 Im z  0

Re z  2 Im z  0 

 2


2

f(!2)

FIGURE 25.10 Strip bounded by vertical lines x = −/2 and x = /2 and the x axis.

FIGURE 25.9

The boundary of S consists of the segment −/2 ≤ x ≤ /2 on the real axis, together with the half-lines x = −/2 and x = /2 for y ≥ 0. We will carry out the strategy of looking at images of the lines bounding S. First, w = u + iv = sin
x cosh
y + i cos
x sinh
y If x = −/2, then w = u + iv = − cosh
y Since 0 ≤ y < on this part of the boundary of S, then cosh
y varies from 1 to . The image of the left vertical boundary of S is therefore the interval
−  −1 on the real axis in the u v plane. If x = /2, a similar analysis shows that the image of the right vertical boundary of S is

1  on the real axis in the u v plane. Finally, if y = 0, then w = sin
x As x varies from −/2 to /2, sin
x varies from −1 to 1. Thus /2 /2 maps to −1 1 in the u v plane. Figure 25.11 shows these results. The boundary of S maps onto the entire real axis in the u v plane. This axis is the boundary of two sets in the w plane, the upper half-plane and the lower half-plane. Choose any convenient z in S, say z = i. Its image is w = sin
i = i sinh
1 lying in the upper half-plane. Therefore the image of S is the upper half-plane. Orientation plays a role in these mappings. Imagine walking along the boundary of S in a counterclockwise sense. This means that we start somewhere up the left boundary x = −/2, walk down this line to the real axis, then turn left and walk along this axis to x = /2, then turn left again and proceed up the right boundary line. Follow the movement of the image point f
z z

w y



Re z 
2 Im z  0


2

v



Re z  2 Im z  0  2

x


1

u 1

w = sin
z maps x = −/2, y ≥ 0, to u ≤ −1; −/2 ≤ x ≤ /2 to −1 ≤ u ≤ 1; and x = /2, y ≥ 0, to u ≥ 1.

FIGURE 25.11

25.1 Functions as Mappings

1061

as z takes this route. As z moves down the line x = −/2, f
z = sin
z begins somewhere to the left of −1 on the real axis in the w plane, and moves toward w = −1. As w turns the first corner and moves along the real axis in the z plane from −/2 to /2, f
z continues through −1 and proceeds along the real axis to w = +1. Finally, z turns upward and moves along the line x = /2, and f
z moves through w = 1 and out the real axis in the w plane. As z traverses the boundary of the strip in a counterclockwise sense (interior of S to the left), f
z traverses the boundary of the upper half plane from left to right in a counterclockwise sense (interior of this plane to the left). In this example, as z moves over the boundary of D in a positive (counterclockwise) sense, f
z moves over the boundary of f
D in a positive sense. In the next section we will discuss mappings that preserve angles and sense of rotation.

SECTION 25.1

PROBLEMS

1. In each of (a) through (e), find the image of the given rectangle under the mapping w = ez . Sketch the rectangle in the z plane and its image in the w plane.

4. Determine the image of the sector /4 ≤  ≤ 5/4 under the mapping w = z2 . Sketch the sector and its image.

(a) 0 ≤ x ≤  0 ≤ y ≤    (b) −1 ≤ x ≤ 1 − ≤ y ≤ 2 2  (c) 0 ≤ x ≤ 1 0 ≤ y ≤ 4 (d) 1 ≤ x ≤ 2 0 ≤ y ≤    (e) −1 ≤ x ≤ 2 − ≤ y ≤ 2 2 2. In each of (a) through (e), find the image of the rectangle under the mapping w = cos
z. Sketch the rectangle and its image in each case.

5. Determine the image of the sector /6 ≤  ≤ /3, under the mapping w = z3 . Sketch the sector and its image.

(a) 0 ≤ x ≤ 1 1 ≤ y ≤ 2  ≤ x ≤  1 ≤ y ≤ 3 (b) 2  (c) 0 ≤ x ≤  ≤ y ≤  2 (d)  ≤ x ≤ 2 1 ≤ y ≤ 2  (e) 0 ≤ x ≤  0 ≤ y ≤ 1 2 3. In each of (a) through (e), find the image of the rectangle under the mapping w = 4 sin
z. Sketch the rectangle and its image in each case.   (a) 0 ≤ x ≤  0 ≤ y ≤ 2 2    ≤ x ≤ 0 ≤ y ≤ (b) 4 2 4  (c) 0 ≤ x ≤ 1 0 ≤ y ≤ 6  3  ≤x≤ 0 ≤ y ≤ (d) 2 2 2 (e) 1 ≤ x ≤ 2 1 ≤ y ≤ 2

6. Show that the mapping w=

  1 1 z+ 2 z

maps the circle z = r onto an ellipse with foci 1 and −1 in the w plane. Sketch a typical circle and its image. 7. Show that the mapping of Problem 6 maps a half-line  = constant onto a hyperbola with foci ±1 in the w plane. Sketch a typical half-line and its image. 8. Show that the mapping w = 1/z maps every straight line to a circle or straight line, and every circle to a circle or straight line. Give an example of a circle that maps to a line, and a line that maps to a circle. 9. Determine the image of the infinite strip 0 ≤ Im
z ≤ 2 under the mapping w = ez . 10. Let D consist of all z in the rectangle having vertices ±i and  ± i, with  a positive number. (a) Determine the image of D under the mapping w = cos
z. Sketch D and its image. (b) Determine the image of D under the mapping w = sin
z. Sketch this image. (c) Determine the image of D under the mapping w = 2z2 . Sketch this image.

1062

25.2

CHAPTER 25

Conformal Mappings

Conformal Mappings Let f  D → D∗ be a mapping.

DEFINITION 25.3

Angle Preserving Mapping

We say that f preserves angles if, for any z0 in D, if two smooth curves in D intersect at z0 , and the angle between these curves at z0 is , then the images of these curves intersect at the same angle  at f
z0 .

This idea is illustrated in Figure 25.12. The images of !1 and !2 are curves f
!1  and f
!2  in D∗ . Suppose !1 and !2 intersect at z0 , and that their tangents have an angle  between them there. We require that the tangents to f
!1  and f
!2  intersect at f
z0  at the same angle. If this condition holds for all smooth curves passing through each point of D, then f is angle preserving on D. z

w

!2 

f(!2)

!1

f(z 0)



z0

FIGURE 25.12

DEFINITION 25.4

f(!1)

Angle-preserving mapping.

Orientation Preserving Mapping

f preserves orientation if a counterclockwise rotation in D is mapped by f to a counterclockwise rotation in D∗ .

This idea is illustrated in Figure 25.13. If L1 and L2 are lines through any point z0 in D, and the sense of rotation from L1 to L2 is counterclockwise, then the sense of rotation from f
L1  to f
L2  through f
z0  in D∗ must also be counterclockwise. Of course f
L1  and f
L2  need not be straight lines, but one can still consider a sense of rotation frrom the tangent to f
L1  to the tangent to f
L2  at f
z0 . By contrast, Figure 25.14 illustrates a mapping that does not preserve orientation. z

L2 z0

L1

FIGURE 25.13

mapping.

w

f (L 2 ) f (z 0)

z

f (L 1)

Orientation-preserving

L2 L1

FIGURE 25.14

mapping.

w f (L 1) f (L 2 )

A nonorientation-preserving

25.2 Conformal Mappings

1063

Preservation of angles and orientation are independent concepts. A mapping may preserve one but not the other. If f  D → D∗ preserves both, we say that f is conformal.

DEFINITION 25.5

Conformal Mapping

f  D → D∗ is a conformal mapping if f is angle preserving and orientation preserving.

The following theorem generates many examples of conformal mappings. THEOREM 25.1

Let f  D → D∗ be a differentiable function defined on a domain D. Suppose f

z = 0 for all z in D. Then f is conformal. Thus, a differentiable function with a nonvanishing derivative on a domain (connected, open set) maps this set in such a way as to preserve both angles and orientation. We will sketch an argument showing why this is true. Let z0 be in D and let ! be a smooth curve in D through z0 . Then f
! is a smooth curve through f
z0  in D∗ (Figure 25.15). If w = f
z and w0 = f
z0 , then w − w0 =

f
z − f
z0 
z − z0  z − z0

Now recall that argument behaves like a logarithm, in the sense that any argument of a product is a sum of arguments of the individual factors. Then   f
z − f
z0  + arg
z − z0  arg
w − w0  = arg (25.1) z − z0 In Figure 25.16,  is the angle between the positive real axis and the line through z and z0 , and is an argument of z − z0 . The angle  between the positive real axis and the line through w and w0 in the w plane is an argument of w − w0 . In the limit as z → z0 , equation (25.1) gives us  = arg f

z0  + . It is here that we use the assumption that f

z0  = 0, because 0 has no argument. If ! ∗ is another smooth curve through z0 , then by the same reasoning, ∗ = arg f

z0  + ∗  Then  − ∗ =  −  ∗  y z

v

w !

z f(!)

z0 D

FIGURE 25.15

f (z 0 )



w w0

z0 D*

θ

FIGURE 25.16

f()

 x

u

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CHAPTER 25

Conformal Mappings

(to within integer multiples of 2). But  − ∗ is the angle between the tangents to ! and ! ∗ at z0 , and  − ∗ is the angle between the tangents to f
! and f
! ∗  at f
z0 . Therefore f preserves angles. The last “equation” also implies that f preserves orientation, since the sense of rotation from ! to ! ∗ is the same as the sense of rotation from f
! to f
! ∗ . A reversal of the sense of rotation would be implied if we had found that  − ∗ = ∗ − . For example, w = sin
z is differentiable, with a nonvanishing derivative, on the strip −/2 < Re
z < /2, and so is a conformal mapping of this strip onto a set in the w plane. A composition of conformal mappings is conformal. Suppose f maps D conformally onto D∗ , and g maps D∗ conformally onto D∗∗ . Then g  f maps D conformally onto D∗∗ (Figure 25.17), because angles and orientation are preserved at each stage of the mapping. We will now consider an important class of conformal mappings.

25.2.1 Linear Fractional Transformations Often we have domains D and D∗ (for example, representing areas of a fluid flow), and we want to produce a conformal mapping of D onto D∗ . This can be a formidable task. Linear fractional transformations are relatively simple conformal mappings that will sometimes serve this purpose.

DEFINITION 25.6

Linear Fractional Transformation

A linear fractional transformation is a function az + b T
z =  cz + d in which a, b, c and d are given complex numbers, and ad − bc = 0.

Other names for this kind of function are Möbius transformation and bilinear transformation. The function is defined except at z = −d/c, which is a simple pole of T . Further, T

z =

ad − bc 
cz + d2

f D*

D

g°f

g

D**

FIGURE 25.17 A composition of conformal mappings is conformal.

25.2 Conformal Mappings

1065

and this is nonzero if z = −d/c. T is therefore a conformal mapping of the plane with the point z = −d/c removed. The condition ad − bc = 0 ensures that T is one-to-one, hence invertible. If we set w =
az + b/
cz + d, then the inverse mapping is z=

dw − b  −cw + a

which is also a linear fractional transformation. We will look at some special kinds of linear fractional transformations.

EXAMPLE 25.6

Let w = T
z = z + b, with b constant. This is called a translation because T shifts z horizontally by Re
b units and vertically by Im
z units. For example, if T
z = z + 2 − i, then T takes z and moves it two units to the right and one unit down (Figure 25.18). We can see this with the following points and their images: 0 → 2 − i

1 → 3 − i

i → 2

4 + 3i → 6 + 2i

EXAMPLE 25.7

Let w = T
z = az, with a a nonzero constant. This is called a rotation/magnitication. To see why, first observe that w = a z  If a > 1, this transformation lengthens a complex number, in the sense of lengthening the line from the origin to z. If a < 1 it shortens this distance. Thus the term magnification. Now write the polar forms z = rei and a = Aei . Then T
z = arei
+  so the transformation adds  to the argument of any nonzero complex number. This rotates the number counterclockwise through an angle . This is the reason for the term rotation. The total effect of the transformation is therefore a scaling and a rotation. As a specific example, consider w =
2 + 2iz This will map i → −2 + 2i 1 → 2 + 2i

and

1 + i → 4i

as shown in Figure 25.19. As Figure 25.20 suggests, in general the image of z is obtained by √ multiplying the magnitude of z by 2 + 2i = 8, and rotating the line from the origin to z counterclockwise through /4 radians. If a = 1, T
z = az is called a pure rotation, since in this case there is no magnification effect, just a rotation through an argument of a.

1066

CHAPTER 25

Conformal Mappings z z z wz 2
i

i

w v

y

4i

1i

2  2i


2  2i

x

u

1

FIGURE 25.18

FIGURE 25.19

y

v

z

w  T(z)

1 z

2  2i z

θ

FIGURE 25.20

z

θ π4

z x

u

The mapping T
z =
2 + 2iz.

1 z

z  1

FIGURE 25.21

Image of z under an inversion.

EXAMPLE 25.8

Let w = T
z = 1/z. This mapping is called an inversion. For z = 0, w =

1 z

and arg
w = arg
1 − arg
z = − arg
z (within integer multiples of 2). This means that we arrive at T
z by moving 1/ z units from the origin along the line from 0 to z, and then reflecting this point across the real axis (Figure 25.21). This maps points inside the unit disk to the exterior of it, and points exterior to the interior, while points on the unit circle remain on the unit circle √ (but get moved around √ this circle, except for 1 and −1). For example, if z =
1 + i/ 2, then 1/z =
1 − i/ 2 (Figure 25.22). We will now show that translations, rotation/magnifications and inversions are the fundamental linear fractional transformations, in the sense that any such mapping can be achieved as a sequence of transformations of these three kinds. To see how to do this, begin with T
z =

az + b  cz + d

T
z =

a b z+ d d

If c = 0, then

is a rotation/magnification followed by a translation: z →

rot/mag

a a b z → z+  d trans d d

25.2 Conformal Mappings

1067

z 1i

z

2

1 z

1
i

z  1

2



1 1i 2

FIGURE 25.22 Image of a point on the unit circle under an inversion.

If c = 0, then T is the result of the following sequence: z → cz → cz + d → rot/mag

→

rot/mag

trans

inv

1 cz + d

a az + b bc − ad 1 bc − ad 1 → + = = T
z c cz + d trans c cz + d c cz + d

This way of breaking a linear fractional transformation into simpler components has two purposes. First, we can analyze general properties of these transformations by analyzing the simpler component transformations. Perhaps more important, we sometimes use this sequence to build conformal mappings between given domains. The following is a fundamental property of linear fractional transformations. We use the term line to mean straight line. THEOREM 25.2

A linear fractional transformation maps any circle to a circle or line, and any line to a circle or line. Because of the preceding discussion, we need to verify this only for translations, rotation/magnifications and inversions. It is obvious geometrically that a translation maps a circle to a circle and a line to a line. Similarly, a rotation/magnification maps a circle to a circle and a line to a line. Now we need to determine the effect of an inversion on a circle or line. Begin with the fact that any circle or line in the plane is the graph of an equation

Proof

A
x2 + y2  + Bx + Cy + R = 0 in which A, B, C, and R are real numbers. This graph is a circle if A = 0 and a line if A = 0 and B and C are not both zero.With z = x + iy , this equation becomes A z2 +

B C
z + z +
z − z + R = 0 2 2i

Now let w = 1/z. The image in the w plane of this locus is the graph of     1 1 1 B 1 C 1 A 2+ + − + + R = 0 2 w w 2i w w w Multiply this equation by ww (the same as w2 ) to get R w2 +

B C
w + w −
w − w + A = 0 2 2i

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CHAPTER 25

Conformal Mappings

In the w plane, this is the equation of a circle if R = 0, and a line if A = 0 and B and C are not both zero. As the proof shows, translations and rotation/magnifications actually map circles to circles and lines to lines, while an inversion may map a circle to a circle or line, and a line to a circle or line.

EXAMPLE 25.9

Let w = T
z = i
z − 2 + 3. This is the sequence z → z − 2 → i
z − 2 → i
z − 2 + 3 = w

(25.2)

a translation by 2 to the left, followed by a counterclockwise rotation by /2 radians (an argument of i is /2), and then a translation by 3 to the right. Since this mapping does not involve an inversion, it maps circles to circles and lines to lines. As a specific example, consider the circle K given by
x − 22 + y2 = 9 with radius 3 and center
2 0. Write this equation as x2 + y2 − 4x − 5 = 0 or z2 − 2
z + z − 5 = 0 Solve w = i
z − 2 + 3 to get z = −i
w − 3 + 2 and substitute into the last equation to get  i
w − 3 + 22 − 2 −i
w − 3 + 2 + −i
w − 3 + 2 − 5 = 0 After some routine manipulation this gives us w2 − 3
w + w = 0 With w = u + iv, this is
u − 32 + v2 = 9 a circle of radius 3 and center
3 0 in the u v− plane. This result could have been predicted geometrically from the sequence of elementary mappings 25.2, shown in stages in Figure 25.23. The sequence first moves the circle 2 units to the left, then (multiplication by i) rotates it by /2 radians (which leaves the center and radius the same), and then translates it 3 units right. The result is a circle of radius 3 about
3 0.

EXAMPLE 25.10

We will examine the effects of an inversion w = 1/z on the vertical line Re
z = a = 0. On this line, z = a + iy and its image consists of points w= It is routine to check that

a 1 y = 2 − 2 i = u + iv 2 a + iy a + y a + y2   1 2 1 + v2 = 2  u− 2a 4a

25.2 Conformal Mappings y


1

Translation x 5 (x
2) 2  y 2  9

2

Rotation
3

(a)

Translation
3

3 (b)

3

1069

v

0

(c)

3

6

u

(d)

FIGURE 25.23

so the image of this vertical line is a circle in the u v plane having center
1/2a 0 and radius 1/2a. In preparation for constructing mappings between given domains, we will show that we can always produce a linear fractional transformation mapping three given points to three given points. THEOREM 25.3

Three Point Theorem

Let z1 , z2 and z3 be three distinct points in the z plane, and w1 , w2 and w3 three distinct points in the w plane. Then there is a linear fractional transformation T of the z plane to the w plane such that T
z1  = w1  T
z2  = w2 Proof

and T
z3  = w3 

Let w = T
z be the solution for w in terms of z and the six given points in the equation
w1 − w
w3 − w2 
z1 − z2 
z3 − z =
z1 − z
z3 − z2 
w1 − w2 
w3 − w

Substitution of z = zj into this equation yields w = wj for j = 1 2 3.

EXAMPLE 25.11

We will find a linear fractional transformation that maps 3 → i 1 − i → 4

and

2 − i → 6 + 2i

Put z1 = 3

z2 = 1 − i

w1 = i

w2 = 4

z3 = 2 − i

and w3 = 6 + 2i

in equation (25.3) to get
i − w
2 + 2i
2 + i
2 − i − z =
3 − z
1
i − 4
6 + 2i − w Solve for w to get w = T
z = Then each T
zj  = wj .


20 + 4iz −
68 + 16i 
6 + 5iz −
22 + 7i

(25.3)

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CHAPTER 25

Conformal Mappings z N: (0, 0, 2) S(x, y) y x

(x, y)

FIGURE 25.24

Stereographic projection identifying the complex sphere with the extended complex plane.

One can show that specification of three points and their images uniquely determines a linear fractional transformation. In the last example, then, T is the only linear fractional transformation mapping the three given points to their given images. When dealing with mappings, it is sometimes convenient to replace the complex plane with the complex sphere. To visualize how this is done, consider the three-dimensional coordinate system in Figure 25.24. A sphere of radius 1 is placed with its south pole at the origin and north pole at
0 0 2. Think of the x y plane as the complex plane. For any
x y in this plane, the line from
0 0 2 to
x y intersects the sphere in exactly one point S
x y. This associates with each point on the sphere, except
0 0 2, a unique point in the complex plane, and conversely. This mapping is called the stereographic projection of the sphere (minus its north pole) onto the plane. This punctured sphere is called the complex sphere. The point
0 0 2 plays the role of a point at infinity. This is motivated by the fact that, as
x y is chosen further from the origin in the x y plane, S
x y moves closer to
0 0 2 on the sphere. The point
0 0 2 is not associated with any complex number, but gives a way of envisioning infinity as a point, something we cannot do in the plane. The extended complex plane (consisting of all complex numbers, together with infinity) is in a one-to-one correspondence with this sphere, including its north pole. To get some feeling for the point at infinity, consider the line y = x in the x y plane. This consists of complex numbers x + xi. If we let x → , the point
1 + ix moves out this line indefinitely far away from the origin. The image of this line on the complex sphere is part of a great circle, and the image point S
x y on the sphere approaches
0 0 2 as x → . This enables us to think of
1 + ix as approaching a specific location we can point to in this limit process, instead of just going further from the origin. In defining a linear fractional transformation, it is sometimes convenient to map one of the three given points in the last theorem to infinity. This can be done by deleting the factors in equation (25.3) involving w3 . THEOREM 25.4

Let z1  z2  z3 be three distinct complex numbers, and w1  w2 distinct complex numbers. Then there is a linear fractional transformation w = T
z that maps T
z1  = w1  T
z2  = w2  Proof

and T
z3  = 

Such a transformation is obtained by solving for w in the equation
w1 − w
z1 − z2 
z3 − z =
z1 − z
w1 − w2 
z3 − z2 

(25.4)

25.2 Conformal Mappings

1071

EXAMPLE 25.12

We will find a linear fractional transformation mapping i → 4i 1 → 3 − i 2 + i →  Solve for w in the equation
4i − w
i − 1
2 + i − z =
i − z
−3 + 5i
1 + i to get w = T
z =


5 − iz − 1 + 3i  −z + 2 + i

Some other properties of linear fractional transformations are pursued in the exercises. We now turn to the problem of constructing conformal mappings between given domains.

SECTION 25.2

PROBLEMS

In each of Problems 1 through 5, find a linear fractional transformation taking the given points to the indicated images. 1. 1 → 1 2 → −i 3 → 1 + i 2. i → i 1 → −i 2 → 0 3. 1 → 1 + i 2i → 3 − i 4 → 4. −5 + 2i → 1 3i → 0 −1 → 5. 6 + i → 2 − i i → 3i 4 → −i In each of Problems 6 through 12, find the image of the given circle or line under the linear fractional transformation. 6. w =

2i  Re
z = −4 z

7. w = 2iz − 4 Re
z = 5 8. w =

z−i 1 1 
z + z +
z − z = 4 iz 2 2i

9. w =

z−1+i  z = 4 2z + 1

10. w = 3z − i z − 4 = 3 11. w =

2z − 5 3  z + z −
z − z − 5 = 0 z+i 2i


1 + 3iz − 2  z − i = 1 12. w = z 13. Prove that the mapping w = z is not conformal. 14. Prove that the composition of two linear fractional transformations is a linear fractional transformation.

15. Prove that every linear fractional transformation has an inverse, and that this inverse is also a linear fractional transformation. (T ∗ is an inverse of T if T  T ∗ and T ∗  T are both the identity mapping, taking each point to itself.) 16. Show that there is no linear fractional transformation mapping the open disk z < 1 onto the set of points bounded by the ellipse u2 /4 + v2 = 1/16. In Problems 17 and 18, the setting is the extended complex plane, which includes the point at infinity. 17. A point z0 is a fixed point of a mapping f if f
z0  = z0 . Suppose f is a linear fractional transformation that is neither a translation nor the identity mapping f
z = z. Prove that f must have either one or two fixed points, but cannot have three. Why does this conclusion fail to hold for translations? How many fixed points can a translation have? 18. Let f be a linear fractional transformation with three fixed points. Prove that f is the identity mapping. In each of Problems 19, 20 and 21, write the linear fractional transformation as the end result of a sequence of mappings, each of which is a translation, rotation/ magnification or inversion. iz − 4 z z−4 20. w = 2z + i 21. w = i
z + 6 − 2 + i z−i 22. w = z+3+i 19. w =

CHAPTER 25

1072

25.3

Conformal Mappings

Construction of Conformal Mappings Between Domains A strategy for solving some kinds of problems (for example, Dirichlet problems) is to find the solution for a “simple” domain (for example, the unit disk), then map this domain conformally to the domain of interest. This mapping may carry the solution for the disk to a solution for the latter domain. Of course, this strategy is predicated on two steps: finding a domain for which we can solve the problem, and being able to map this domain to the domain for which we want the solution. We will now discuss the latter problem. Although in practice it may be an imposing task to find a conformal mapping between given domains, the following result assures us that such a mapping exists, with one exception.

THEOREM 25.5

Riemann Mapping Theorem

Let D∗ be a domain in the w plane, and assume that D∗ is not the entire w plane. Then there exists a one-to-one conformal mapping of the unit disk z < 1 onto D∗ . This powerful result implies the existence of a conformal mapping between given domains. Suppose we want to map D onto D∗ (neither of which is the entire plane). Insert a third plane, the " plane, between the z plane and the w plane, as in Figure 25.25. By Riemann’s theorem, there is a one-to-one conformal mapping g of the unit disk " < 1 onto D∗ . Similarly, there is a one-to-one conformal mapping f of " < 1 onto D. Then g  f −1 is a one-to-one conformal mapping of D onto D∗ . In theory, then, two domains, neither of which is the entire plane, can be mapped conformally in one-to-one fashion onto one another. This does not, however, make such mappings easy to find. In attempting to find such mappings, the following observation is useful. A conformal mapping of a domain D onto a domain D∗ will map the boundary of D to the boundary of D∗ . We use this as follows. Suppose D is bounded by a path C (not necessarily closed) which separates the z plane into two domains, D and . These are called complementary domains. Similarly, suppose D∗ is bounded by a path C ∗ which separates the w plane into complementary domains D∗ and ∗ (Figure 25.26). Try to find a conformal mapping f that sends points of C to points of C ∗ . This may be easier than trying to find a mapping of the entire domain. This mapping will then send D to either D∗ or to ∗ . To see which, choose a point z0 in D and see whether f
z0  is in D∗ or ∗ . If f
z0  is in D∗ (Figure 25.27(a)), then f  D → D∗ and we have our conformal mapping. If f
z0  is in ∗ as in Figure 25.27(b), then f  D → ∗ . This is not yet the mapping we want, but sometimes we can take another step and use f to manufacture a conformal mapping from D to D∗ . We will now construct some conformal mappings, beginning with very simple ones and building up to more difficult problems.



z D

f
1 f

g  f
1 g

i

w D*

D 1

Mapping D onto D∗ through the unit disk.

FIGURE 25.25

D

z

C FIGURE 25.26

D*

w D*

C*

25.3 Construction of Conformal Mappings Between Domains

z

z

w

D

D*

D z0

f : D → D* D

f (z 0 )

C

C

C*

FIGURE 25.27(a)

D*

D*

D

D*

w

f : D → D*

z0

f (z 0 )

1073

C*

FIGURE 25.27(b)

EXAMPLE 25.13

Suppose we want to map the unit disk D  z < 1 conformally onto the disk D∗  w < 3. Clearly a magnification w = f
z = 3z will do this, because all we have to do is expand the unit disk to a disk of radius 3 (Figure 25.28). Notice that this mapping carries the boundary of D onto the boundary of D∗ .

EXAMPLE 25.14

Map the unit disk D  z < 1 conformally onto the domain w > 3. Here we are mapping D to the complementary domain of the preceding example. We already know that f
z = 3z maps D conformally onto w < 3. Combine this map with an inversion, letting 3 g
z = f
1/z =  z This maps z < 1 to w > 3 (Figure 25.29). Again, the boundary of the unit disk maps to the boundary of w > 3, which is the circle of radius 3 about the origin in the w plane.

EXAMPLE 25.15

We will map the unit disk D  z < 1 onto the disk D∗  w − i < 3, of radius 3 and centered at i in the w plane. Figure 25.30 suggests one way to construct this map. We want to expand the unit disk’s radius by a factor of 3, then translate the resulting disk up one unit. Thus, map in steps: z → 3z → 3z + i w z

y

w  3z

z

w v 3i

x

FIGURE 25.28

onto w < 3.

3 z

3i

i

i 1

w

3

u

Mapping of z < 1

1

3

FIGURE 25.29 Mapping of z < 1

onto w > 3.

1074

CHAPTER 25

Conformal Mappings

w  "  i  3z  i z

"  3z

w

"

4i

w  "  i

3i

y

z

i

i 0
i

3
2i FIGURE 25.30

Mapping of z < 1 onto w − i < 3.

w v i

Re(z)  0 x

w 1
1

1

FIGURE 25.31

a magnification followed by a translation. The mapping is w = f
z = 3z + i. This maps the unit circle z = 1 to the circle w − i = 3 because w − i = 3z = 3 z = 3 Further, the origin in the z plane (center of D) maps to i in the w− plane, and i is the center of D∗ , so f  D → D∗ .

EXAMPLE 25.16

Suppose we want to map the right half-plane D  Re
z > 0 onto the unit disk D∗  w < 1. The domains are shown in Figure 25.31. The boundary of D is the imaginary axis Re
z = 0. We will map this to the boundary of the unit disk w = 1. To do this, pick three points on Re
z = 0, and three on w = 1, and use these to define a linear fractional transformation. There is a subtlety, however. To maintain positive orientation (counterclockwise on closed curves), choose three points in succession down the imaginary axis, so a person walking along these points sees the right half-plane on the left. Map these to points in order counterclockwise around w = 1 For example, we can choose z1 = i

z2 = 0

and z3 = −i

on the imaginary axis in the z plane, and map these in order to w1 = 1

w2 = i

w3 = −1

From equation (25.3), we have
1 − w
−1 − i
i
−i − z =
i − z
−i
1 − i
−1 − w Solve for w:



 z−1 w = T
z = −i  z+1

This conformal mapping must take the right half-plane to the interior or exterior of the unit disk in the w plane. To see which it is, pick a point in Re
z > 0, say z = 1. Since T
1 = 0 is in D∗ , T maps the right half-plane to the unit disk w < 1 as we want.

EXAMPLE 25.17

Suppose we want to map the right half-plane to the exterior of the unit disk, the domain w > 1. We have T  Re
z > 0 → w < 1 from the preceding example. If we follow this map (taking

u

25.3 Construction of Conformal Mappings Between Domains

1075

Re
z > 0 onto the unit disk) with an inversion (taking the unit disk to the exterior of the unit disk), we will have the map we want. Thus, with T
z as in the last example, let   z+1 1 =i  f
z = T
z z−1 As a check, 1 + i is in the right half-plane, and   2+i f
1 + i = i = 2+i i is exterior to the unit disk in the w plane.

EXAMPLE 25.18

We will map the right half-plane Re
z > 0 conformally onto the disk w − i < 3. We can do this as a composition of mappings we have already constructed. Put an intermediate " plane between the z and w planes (Figure 25.32). From Example 25.16, map Re
z > 0 onto the unit disk " < 1 by   z−1 " = f
z = −i  z+1 Now use the mapping of Example 25.15 to send the unit disk " < 1 onto the disk w − i < 3: w = g
" = 3" + i The composition g  f is a conformal mapping of Re
z > 0 onto w − i < 3:   z−1 2i
−z + 2 +i =  w =
g  f
z = g
f
z = 3f
z + i = 3 −i z+1 z+1

EXAMPLE 25.19

We will map the infinite strip S  −/2 < Im
z < /2 onto the unit disk w < 1. Recall from Example 25.4 that the exponential function maps horizontal lines to halflines from the origin. The boundary of S consists of two horizontal lines, Im
z = −/2 and Im
z = /2. On the lower boundary line, z = x − i/2, so ez = ex e−i/2 = −iex  which varies over the negative imaginary axis as x takes on all real values. On the upper boundary of S, z = x + i/2, and ez = iex varies over the positive part of the imaginary axis as x varies over the real line. The imaginary axis forms the boundary of the right half-plane Re
z > 0, as well as of the left half-plane Re
w < 0 in the w plane. The mapping w = ez must map S to one of these complementary domains. However, this mapping sends 0 to 1, in the right half-plane, so the mapping w = f
z = ez maps S to the right half-plane. We want to map S onto the unit disk. But now we know a mapping of S onto the right half-plane, and we also know a mapping of the right half-plane onto the unit disk. All we have to do is put these together.

1076

CHAPTER 25

Conformal Mappings

z

gf

w 4i

g

i

"

f

i

z

w 1

x

1

u Re(z" )  0

– i


2i FIGURE 25.32

w v

" i 2

2

 Im
z < 2 → Re
" > 0 → w < 1 yields a mapping of  Im
z < 2 onto w < 1.

Mapping of Re
z > 0

FIGURE 25.33

onto w − i < 3.

In Figure 25.33, put a " plane between the z and w planes. Map " = f
z = ez taking S onto the right half-plane Re
" > 0. Next map   " −1 w = g
" = −i  " +1 taking the right half-plane Re
" > 0 onto the unit disk w < 1. Therefore, the function w =
g  f
z = g
f
z   z e −1 = g
ez  = −i z e +1 is a conformal mapping of S onto w < 1. In terms of hyperbolic functions, this mapping can be written w = −i tanh
z/2 The conformal mapping of the last example is not a linear fractional transformation. These are convenient to use whenever possible. However, even when we know from the Riemann mapping theorem that a conformal mapping exists between two domains, we are not guaranteed that we can always find such a mapping in the form of a linear fractional transformation.

EXAMPLE 25.20

We will map the disk z < 2 onto the domain D∗  u + v > 0 in the u v plane. These domains are shown in Figure 25.34. Consider mappings we already have that relate to this problem. First, we can map z < 2 to " < 1 by a simple magnification (multiply by 21 ). But we also know a mapping from the unit disk to the right half-plane. Finally, we can obtain D∗ from the right half-plane by a

y

z

w v D*

D

u

x z2

FIGURE 25.34

u  v 0

25.3 Construction of Conformal Mappings Between Domains z

" 1 w

z 2

w



"

1077

 4

Re (z )  0

u  v 0

FIGURE 25.35

counterclockwise rotation through /4 radians, an effect achieved by multiplying by ei/4 . This suggests the strategy of constructing the mapping we want in the stages shown in Figure 25.35: z < 2 → " < 1 → Re
 > 0 → u + v > 0 The first step is achieved by 1 " = z 2 Next, use the inverse of the mapping from Example 25.16 and name the variables " and  to get =

1 + i"  1 − i"

This maps " < 1 → Re
 > 0 Finally, perform the rotation: w = ei/4  In sum,

 w=e

i/4

=e

i/4

1 + i" 1 − i"



 =e

i/4

1 + i
z/2 1 − i
z/2

 =

2 + iz i/4 e  2 − iz

This maps the disk z < 2 conformally onto the half-plane u + v > 0. For example, 0 is in the disk, and √ 2 i/4
1 + i w
0 = e = 2 is in u + v > 0. We will briefly discuss the Schwarz–Christoffel transformation, which may be used when a domain has a polygon as a boundary.

25.3.1

Schwarz–Christoffel Transformation

Suppose we want a conformal mapping of the upper half-plane to the interior  of a polygon P, which could be a triangle, rectangle, pentagon or other polygon. A linear fractional transformation will not do this. However, the Schwarz–Christoffel transformation was constructed just for this purpose. Let P have vertices w1      wn in the w plane (Figure 25.36). Let the exterior angles of P be 1      n . We claim that there are constants z0 , a and b, with Im
z0  > 0, and real numbers x1      xn such that the function
z f
z = a
 − x1 −1
 − x2 −2 · · ·
 − xn −n d + b (25.5) z0

1078

CHAPTER 25

Conformal Mappings z

w v

 n

 1 w2

w1

y

1

w v

wn

 2 xn

x1 x2

x

wn

2 w3

u

3

u FIGURE 25.36

w1 n

w2

FIGURE 25.37

is a conformal mapping of onto . This integral is taken over any path in from z0 to z in . The factors
 − xj −j are defined using the complex logarithm obtained by taking the argument lying in 0 2. Any function of the form of equation (25.5) is called a Schwarz–Christoffel transformation. To see the idea behind this function, suppose each xj < xj+1 . If z is in , let g
z = a
z − x1 −1
z − x2 −2 · · ·
z − xn −n  Then f

z = g
z and arg f

z = arg
z − 1 arg
z − x1  − · · · − n arg
z − xn  As we saw in the discussion of Theorem 25.1, arg f

z is the number of radians by which the mapping f rotates tangent lines, if f

z = 0. Now imagine z moving from left to right along the real axis (Figure 25.37), which is the boundary of . On
−  x1 , f
z moves along a straight line (no change in the angle). As z passes over x1 , however, arg f

z changes by 1 . This angle remains fixed as z moves from x1 toward x2 . As z passes over x2 , arg f

z changes by 2 , then remains at this value until z reaches x3 , where arg f

z changes by 3 , and so on. Thus arg f

z remains constant on intervals
xj−1  xj  and increases by j  as z passes over xj . The net result is that the real axis is mapped to a polygon P ∗ having exterior angles 1      n . These numbers are actually determined by 1      n−1 , since n 

j  = 2

j=1

P ∗ has the same exterior angles as P but need not be the same as P because of its location and size. We may have to rotate, translate and/or magnify P ∗ to obtain P. These effects are achieved by choosing x1      xn to make P ∗ similar to P, and then choosing a (rotation/magnification) and b (translation) to obtain P. If we choose zn = , then z1      zn−1  are mapped to the vertices of P. In this case the Schwarz–Christoffel transformation is
z f
z = a
 − x1 −1
 − x2 −2 · · ·
 − xn −n−1 d + b (25.6) z0

It can be shown that any conformal mapping of onto a polygon must have the form of a Schwarz–Christoffel transformation. In practice a Schwarz–Christoffel transformation can be difficult or impossible to determine in closed form because of the integration.

25.3 Construction of Conformal Mappings Between Domains

1079

EXAMPLE 25.21

We will map the upper half-plane onto a rectangle. Choose x1 = 0 x2 = 1 and x3 as any real number greater than 1. The Schwarz–Christoffel transformation of equation (25.6) has the form
z 1 f
z = a d + b  z0 
 − 1
 − x3  with a and b chosen to fit the dimensions of the rectangle and its orientation with respect to the  axes. The radical appears because the exterior angles of a rectangle are all equal to /2, so 4j=1 j = 4k = 2. This integral is an elliptic integral and cannot be evaluated in closed form.

EXAMPLE 25.22

We will map onto the strip S  Im
w > 0 −c < Re
w < c in the w plane. Here c is a positive constant.

and the strip S are shown in Figure 25.38. To use the Schwarz–Christoffel transformation, we must think of S as a polygon with vertices −c, c and . Choose x1 = −1 to map to −c and x2 = 1 to map to c. Map to . The exterior angles of the strip are /2 and /2, so 1 = 2 = 21 . The transformation has the form
z w = f
z = a
 + 1−1/2
 − 1−1/2 d + b z0

We will choose z0 = 0 and b = 0. Write
 − 1−1/2 = −
1 − −1/2 = −i
1 − −1/2  With −ai = A, we have w = f
z = A




z 0

1 d
1 −  2 1/2

This integral is reminiscent of the real integral representation of the inverse sine function. Indeed, we can write w = A sin−1
z by which we mean that z = sin

w

 A To choose A so that −1 maps to −c and 1 to c, we need c = 1 sin A

y

z

w v x

FIGURE 25.38

Re w 
c
c

Re w  c c

u

CHAPTER 25

1080

Conformal Mappings

Thus choose c/A = /2, or A=

2c  

The mapping is w=

2c −1 sin
z 

If we choose c = /2, this mapping is just w = sin−1
z, mapping onto the strip Im
w > 0, −/2 < Re
w < /2. This is consistent with the result of Example 25.5.

SECTION 25.3

PROBLEMS

In each of Problems 1 through 6, find a linear fractional transformation mapping the first domain onto the second. 1. z < 3 onto w − 1 + i < 6 2. z < 3 onto w − 1 + i > 6

maps the upper half plane onto the rectangle with  verticesl   0, c, c + ic and ic, where c = 1 1 3 $ $ /$ . Here $ is the gamma function. 2 4 4 10. Define the cross ratio of z1  z2  z3 and z4 to be the image of z1 under the linear fractional transformation that maps z2 → 1, z3 → 0, z4 → . Denote this cross ratio as z1  z2  z3  z4 . Suppose T is any linear fractional transformation. Show that T preserves the cross ratio. That is,

3. z + 2i < 1 onto w − 3 > 2 4. Re
z > 1 onto Im
w > −1 5. Re
z < 0 onto w < 4 6. Im
z > −4 onto w − i > 2 7. Find a conformal mapping of the upper half-plane Im
z > 0 onto the wedge 0 < arg
w < /3.

z1  z2  z3  z4  = Tz1  Tz2  Tz3  Tz4 

8. Let w = Log
z, in which the logarithm is given a unique value for each nonzero z by rectricting the argument of z to lie in 0 2. Show that this mapping takes Im
z > 0 onto the strip 0 < Im
w < .

11. Prove that z1  z2  z3  z4  is the image of z1 under the linear fractional transformation defined by z − z4 z − z2  w = 1− 3 z3 − z2 z − z4

9. Show that the Schwarz–Christoffel transformation
z f
z = 2i
 + 1−1/2
 − 1−1/2  −1/2 d

12. Prove that z1  z2  z3  z4  is real if and only if the z
j s are on the same circle or straight line.

0

25.4

Harmonic Functions and the Dirichlet Problem Given a set D of points in the plane, let D denote the boundary of D. A Dirichlet problem for D is to find a solution of Laplace’s equation

2 u 2 u + =0

x2 y2 for
x y in D, satisfying the boundary condition u
x y = f
x y for
x y in D Here f is a given function, usually assumed to be continuous on the boundary of D.

25.4 Harmonic Functions and the Dirichlet Problem

1081

A function satisfying Laplace’s equation in a set is said to be harmonic on that set. Thus the Dirichlet problem for a set is to find a function that is harmonic on that set, and satisfies given data on the boundary of the set. Chapter 19 is devoted to solutions of Dirichlet problems using methods from real analysis. Our purpose here is to apply complex function methods to Dirichlet problems. The connection between a Dirichlet problem and complex function theory is given by the following.

THEOREM 25.6

Let D be an open set in the plane, and let f
z = u
x y + iv
x y be differentiable on D. Then u and v are harmonic on D. That is, the real and imaginary parts of a differentiable complex function are harmonic. Proof

By the Cauchy–Riemann equations,    

2 u 2 u

u

u + = +

x2 y2

x x

y y    

2 v

v

v

2 v − = 0 = + − =

x y

y

x

x y y x

Therefore u is harmonic on D. The proof that v is harmonic is similar. Conversely, given a harmonic function u, there is a harmonic function v so that f
z = u
x y + iv
x y is differentiable. Such a v is called a harmonic conjugate for u. THEOREM 25.7

Let u be harmonic on a domain D. Then, for some v, u
x y + iv
x y defines a differentiable complex function for z = x + iy in D. Proof

Let g
z =

u

u −i

x

y

for
x y in D. Using the Cauchy–Riemann equations and Theorem 21.6, we find that g is differentiable on D. Then, for some function G, G

z = g
z for z in D. Write G
z = U
x y + iV
x y Now G

z =

U

u

u

U −i = g
z = −i 

x

y

x

y

Therefore

U

u =

x

x

and

U

u =

y

y

on D. Then, for some constant K, U
x y = u
x y + K

1082

CHAPTER 25

Conformal Mappings

Let f
z = G
z − K. Then f is differentiable at all points of D. Further, f
z = G
z − K = U
x y + iV
x y − K = u
x y + iv
x y We may therefore choose v
x y = V
x y, proving the theorem. Given a harmonic function u, we will not be interested in actually producing a harmonic conjugate v. However, we will exploit the fact that there exists such a function to produce a differentiable complex function f = u + iv, given harmonic u. This enables us to apply complex function methods to Dirichlet problems. As a preliminary, we will derive two important properties of harmonic functions. THEOREM 25.8

Mean Value Property

Let u be harmonic on a domain D. Let
x0  y0  be any point of D, and let C be a circle of radius r centered at
x0  y0 , contained in D and enclosing only points of D. Then 1
2 u
x0  y0  = u
x0 + r cos
 y0 + r sin
d 2 0 As  varies from 0 to 2,
x0 + r cos
 y0 + r sin
 moves once about the circle of radius r about
x0  y0 . The conclusion of the theorem is called the mean value property because it states that the value of a harmonic function at the center of any circle in the domain is the average of its values on the circle. For some v, f = u + iv is differentiable on D. Let z0 = x0 + iy0 . By Cauchy’s integral formula, 1 - f
z u
x0  y0  + iv
x0  y0  = f
z0  = dz 2i C z − z0 1
2 f
z0 + rei  i = ire d 2i 0 rei 1
2 = u
x0 + r cos
 y0 + r sin
d 2 0 i
2 + v
x0 + r cos
 y0 + r sin
d 2 0

Proof

By taking the real and imaginary part of both sides of this equation, we get the conclusion of the theorem. If D is a bounded domain, then the set D consisting of D, together with all boundary points of D, is called the closure of D. D is a closed and bounded, hence compact set. If u
x y is continuous on D, then u
x y must achieve a maximum value on D. If u is also harmonic on D, we claim that this maximum must occur at a boundary point of D. This is reminiscent of the maximum modulus theorem, from which it follows. THEOREM 25.9

Let D be a bounded domain. Suppose u is continuous on D and harmonic on D. Then u
x y achieves its maximum value on D at a boundary point of D. Proof

First produce v so that f = u + iv is differentiable on D. Define g
z = ef
z

25.4 Harmonic Functions and the Dirichlet Problem

1083

for all z in D. Then g is differentiable on D. By the maximum modulus theorem, g
z achieves its maximum at a boundary point of D. But   g
z = eu
xy+iv
xy  = eu
xy  Since eu
xy is a strictly increasing real function, eu
xy and u
x y must achieve their maximum values at the same point. Therefore u
x y must achieve its maximum at a boundary point of D. For example, u
x y = x2 −y2 is harmonic on the open unit disk x2 +y2 < 1, and continuous on its closure x2 +y2 ≤ 1. This function must therefore achieve its maximum value for x2 +y2 ≤ 1 at a boundary point of this disk, namely at a point for which x2 + y2 = 1. We find that this maximum value is 1, occurring at
1 0 and at
−1 0.

25.4.1

Solution of Dirichlet Problems by Conformal Mapping

We want to use conformal mappings to solve Dirichlet problems. The strategy is to first solve the Dirichlet problem for a disk. Once we have this, we can attempt to solve a Dirichlet problem for another domain D by constructing a conformal mapping between the unit disk and D, and applying this mapping to the solution we have for the disk. In Section 19.3 we derived Poisson’s integral formula for the solution of the Dirichlet problem for a disk, using Fourier methods. We will use complex function methods to obtain a form of this solution that is particularly suited to using with conformal mappings. ,  z < 1 and assumes given values We want a function u that is harmonic on the disk D u
x y = g
x y on the boundary circle. Suppose u is harmonic on the slightly larger disk z < 1 + . If v is a harmonic conjugate of u, then f = u + iv is differentiable on this disk. We may, by adding a constant if necessary, choose v so that v
0 0 = 0. Expand f in a Maclaurin series f
z =



an zn 

(25.7)

n=0

Then

1 f
z + f
z 2 1 1
an zn + an zn  = a0 +
a zn + an zn   = 2 n=0 2 n=1 n

u
x y = Re
f
x + iy =

Now let " be on the unit circle !. Then "2 = "" = 1, so " = 1/" and the series is u
" = a0 +

1
a " n + an " −n   2 n=1 n

Multiply this equation by " m /2i and integrate over !. Within the open disk of convergence, the series and the integral can be interchanged. We get   1 a0 - m 1 1  m n+m −n+m u
"" d" = " d" + " d" + an " d"  (25.8) a 2i ! 2i ! 2 2i n=1 n ! ! Recall that !

 " d" = k

0 2i

if k = −1  if k = −1

1084

CHAPTER 25

Conformal Mappings

Therefore, if m = −1 in equation (25.8), we have 1 1 u
" d" = a0  2i ! " If m = −n − 1 with n = 1 2 3   , we obtain 1 1 u
"" −n−1 = an  2i ! 2 Substitute these coefficients into equation (25.7) to get     1 1 1 n −n−1 f
z = d"zn an z = u
" d" + u
"" 2i ! " n=0 n=1 i !    n  z 1 u
" d" = 1+2 2i ! " n=1 " Since z < 1 and " = 1, then z/" < 1 and the geometric series in this equation converges:  n  z z z/" =  = " 1 − z/" " − z n=1 Then



 1 " +z 1 1 2z 1 f
z = d" = d" u
" 1 + u
" 2i ! " −z " 2i ! " −z "

If u
" = g
", given values for u on the boundary of the unit disk, then, for z < 1,     1 " +z 1 u
x y = Re f
z = Re d"  g
" 2i ! " −z "

(25.9)

This is an integral formula for the solution of the Dirichlet problem for the unit disk. We leave it as an exercise for the student to retrieve the Poisson integral formula from this expression by putting z = rei and " = ei . Equation (25.9) is well suited to solving certain Dirichlet problems by conformal mappings. , where D , is the Suppose we know a differentiable, one-to-one conformal mapping T  D → D, unit disk w < 1 in the w plane. Assume that T maps C, the boundary of D, onto the unit , and that T −1 is also a differentiable conformal mapping. circle , C bounding D, To help follow the discussion, we will use " for an arbitrary point of , C,  for a point on C, and
, x, y for a point in the w plane (Figure 25.39). Now consider a Dirichlet problem for D:

2 u 2 u + =0

x2 y2

for
x y in D

u
x y = g
x y for
x y in C

y

z

w v

 D x C

FIGURE 25.39


w D:

T

"

T
1

C




y) (x, u

25.4 Harmonic Functions and the Dirichlet Problem

1085

If w = T
z, then z = T −1
w and we define , g
w = g
T −1
w = g
z In the w plane, we now have a Dirichlet problem for the unit disk:

2, u 2, u + 2 =0 2

x

y

, for
, x, y in D

, u
, x, y = , g
, x, y for
, x, y on , C From equation (25.9) the solution of this problem for the unit disk is the real part of   " +w 1 1 , , f
w = d" g
" 2i ,C " −w " Finally, recall that T maps C onto , C, let " = T
 for  on C to obtain u
x y = Re f
z = Re



 1
1
T
 + T
z , T
d  g
T
 2i C T
 − T
z T


Since , g
T
 = g
T −1
T
 = g
, we have the solution u
x y = Re f
z = Re

  1
T
 + T
z T

 d  g
 2i C T
 − T
z T


(25.10)

This solves the Dirichlet problem for the original domain D. To illustrate this technique, we will solve the Dirichlet problem for the right half plane:

2 u 2 u + =0

x2 y2

for x > 0 − < x < 

u
0 y = g
y for − < y <  We need a conformal mapping from the right half plane to the unit disk. There are many such mappings. From Example 25.16, we can use 

 z−1 w = T
z = −i  z+1 Compute T

z =

−2i 
z + 12

From equation (25.10), the solution is the real part of   1
−i
 − 1/
 + 1 − i
z − 1/
z + 1 1 −2i f
z = u
 d 2i C −i
 − 1/
 + 1 + i
z − 1/
z + 1 −i
 − 1/
 + 1
 + 12   1 1
z − 1 = u
 d i C  − z 2 − 1

1086

CHAPTER 25

Conformal Mappings

The boundary C of the right half plane is the imaginary axis, and is not a closed curve. Parametrize C as  =
0 t = it, with t varying from − to for positive orientation on D (as we walk down this axis, D is over our left shoulder). We get    1
− −1 itz − 1 idt f
z = u
0 t i it − z 1 + t2    itz − 1 1
1 = u
0 t dt  − it − z 1 + t2 The solution is the real part of this integral. Now t, u
0  and 1/
1+t2  are real, so concentrate on the term containing i and z = x + iy: itz − 1 itx − ty − 1 itx − ty − 1 −it − x + iy = = it − z it − x − iy it − x − iy −it − x + iy =

tx
t − y − itx2 + ity
t − y + txy + i
t − y + x  x2 +
t − y2

The real part of this expression is x
1 + t2  x2 +
t − y2 Therefore u
x y = Re f
z = =



1
1 x
1 + t2  u
0 t 2 dt 2  − x +
t − y 1 + t2

1
x g
t 2 dt  − x +
t − y2

This is an integral formula for the solution of the Dirichlet problem for the right half-plane.

SECTION 25.4

PROBLEMS

1. Using complex function methods, write an integral solution for the Dirichlet problem for the upper half-plane Im
z > 0. 2. Using complex function methods, write an integral solution for the Dirichlet problem for the right quarterplane Re
z > 0, Im
z > 0, if the boundary conditions are u
x 0 = f
x and u
0 y = 0. 3. Write an integral solution for the Dirichlet problem for the disk z − z0  < R. 4. Write a formula for the solution of the Dirichlet problem for the right half-plane if the boundary condition is given by  1 for − 1 ≤ y ≤ 1 u
0 y = 0 for y > 1 5. Write a formula for the solution of the Dirichlet problem for a unit disk if the boundary condition is given by u
x y = x − y for
x y on the unit circle.

6. Write a formula for the lem for the unit disk given by  1 i u
e  = 0

solution of the Dirichlet probif the boundary condition is for 0 ≤  ≤ /4 for /4 <  < 2

7. Write a formula for the solution of the Dirichlet problem for the strip −1 < Im
z < 1, Re
z > 0, if the boundary condition is given by u
x 1 = u
x −1 = 0 u
0 y = 1 − y

for 0 < x < for − 1 ≤ y ≤ 1

8. Write an integral formula for the solution of the Dirichlet problem for the strip −1 < Re
z < 1, Im
z > 0 if the boundary condition is given by u
x 0 = 1 u
−1 y = u
1 y = e

−y

for − 1 < x < 1 for 0 < y < 

25.5 Complex Function Models of Plane Fluid Flow

25.5

1087

Complex Function Models of Plane Fluid Flow We will discuss how complex functions and integration are used in modeling and analyzing the flow of fluids. Consider an incompressible fluid, such as water under normal conditions. Assume that we are given a velocity field V
x y in the plane. By assuming that the flow depends only on two variables, we are taking the flow to be the same in all planes parallel to the complex plane. Such a flow is called plane-parallel. This velocity vector is also assumed to be independent of time, a circumstance described by saying that the flow is stationary. Write V
x y = u
x yi + v
x yj Since we can identify vectors and complex numbers, we will by a mild abuse of notation write the velocity vector as a complex function V
z = V
x + iy = u
x y + iv
x y Given V
z, think of the complex plane as divided into two sets. The first is the domain D on which V is defined. The complement of D consists of all complex numbers not in D. Think of the complement as comprising channels confining the fluid to D, or as barriers through which the fluid cannot flow. This enables us to model fluid flow through a variety of configurations and around barriers of various shapes. Suppose ! is a closed path in D. From vector analysis, if we parametrize ! by x = x
s y = y
s, with s arc length along the path, then the vector x

si + y

sj is a unit tangent vector to !, and   dx dy i + j ds = udx + vdy
ui + vj · ds ds This is the dot product of the velocity with the tangent to a path trajectory, leading us to interpret udx + vdy !

as a measure of the velocity of the fluid along !. The value of this integral is called the circulation of the fluid along !. The vector −y

si + x

sj is a unit normal vector to !, being perpendicular to the tangent vector (Figure 25.40). Therefore   dy dx −
ui + vj · − i + j ds = −vdx + udy dx ds ! ! is the negative of the integral of the normal component of the velocity along the path. When this integral is not zero, it is called the flux of the fluid across the path. This gives a measure

N

 2

γ FIGURE 25.40

T

1088

CHAPTER 25

Conformal Mappings

of fluid flowing across ! out from the region bounded by !. When this flux is zero for every closed path in the domain of the fluid, the fluid is called solenoidal. A point z0 =
x0  y0  is a vortex of the fluid if the circulation has a nonzero constant value on every closed path about z0 , in the interior of some punctured disk 0 < z − z0  < r. The constant/ value of the circulation is the strength of the vortex. If ! −vdx + udy has the same positive value k for all closed paths about z0 in some punctured disk about z0 , then we call z0 a source of strength k; if k is negative, z0 is a sink of strength k. The connection between the velocity field of a fluid and complex functions is provided by the following. THEOREM 25.10

Let u and v be continuous with continuous first and second partial derivatives in a simply connected domain D. Suppose ui + vj is irrotational and solenoidal in D. Then u and −v satisfy the Cauchy-Riemann equations in D, and f
z = u
x y − iv
x y is a differentiable complex function on D. Conversely, if u and −v satisfy the Cauchy-Riemann equations on D, then ui + vj defines an irrotational, solenoidal flow on D. Let ! be any closed path in D. If M is the interior of !, then every point in M is also in D by the assumption that D is simply connected. By Green’s theorem,

 v u  − udx + vdy = dA = 0

x y ! M

Proof

because the flow is irrotational. But the flow is also solenoidal, so, again by Green’s theorem,

 u v  + dA = 0 −vdx + udy =

x y ! M Because M can be any set of points in D bounded by a closed path, the integrands in both of these double integrals must be zero throughout D, so

u =
−v

x y

and

u

= −
−v

y

x

By Theorem 21.6, f
z = u
x y − iv
x y is differentiable on D. The converse follows by a similar argument. Theorem 25.9 provides some insight into irrotational, solenoidal flows. If the flow is irrotational, then  

v u − curl
ui + vj = k = 0

x y as shown in the proof of the theorem. This curl is a vector normal to the plane of the flow. From the discussion of Section 12.5.2, the curl of ui + vj is twice the angular velocity of the particle of fluid at
x y. The fact that this curl is zero for an irrotational flow means that the fluid particles may experience translations and distortions in their motion, but no rotation. There is no swirling effect in the fluid. If the flow is solenoidal, then div
ui + vj =

u v + = 0

x y

A further connection between flows and complex functions is provided by the following.

25.5 Complex Function Models of Plane Fluid Flow

1089

THEOREM 25.11

Let f be a differentiable function defined on a domain D. Then f

z is an irrotational, solenoidal flow on D. Conversely, if V = ui + vj is an irrotational, solenoidal vector field on a simply connected domain D, then there is a differentiable complex function f defined on D such that f

z = V . Further, if f
z = 
x y + i
x y, then





 = u = v = −v

x

y

x

and

 = u

y

We leave a proof of this result to the student. In view of the fact that f

z is the velocity of the flow, we call f a complex potential of the flow. Theorem 25.11 implies that any differentiable function f
z = 
x y + i
x y defined on a simply connected domain determines an irrotational, solenoidal flow f

z =



 +i = u
x y − iv
x y = u
x y + iv
x y

x

x

We call  the velocity potential of the flow, and curves 
x y = k are called equipotential curves. The function  is called the stream function of the flow, and curves 
x y = c are called streamlines. We may think of w = f
z as a conformal mapping wherever f

z = 0. A point at which
f
z = 0 is called a stagnation point of the flow. Thinking of f as a mapping, in the w plane we have w = f
z = 
x y + i
x y =  + i Equipotential curves 
x y = k map under f to vertical lines  = k, and streamlines 
x y = c map to horizontal lines  = c. Since these sets of vertical and horizontal lines are mutually orthogonal in the w plane, the streamlines and equipotential curves in the z plane also form orthogonal families. Every streamline is orthogonal to each equipotential curve at any point where they intersect. This conclusion fails at a stagnation point, where the mapping may not be conformal. Along an equipotential curve 
x y = k, d =



 dx + dy = udx + vdy = 0

x

y

Now ui + vj is the velocity of the flow at
x y and x

si + y

sj is a unit tangent to the equipotential curve through
x y. Since the dot product of these two vectors is zero from the fact that d = 0 along the equipotential curve, we conclude that the velocity is orthogonal to the equipotential curve through
x y, provided that
x y is not a stagnation point. Similarly, along a streamline 
x y = c, d =



 dx + dy = −vdx + udy = 0

x

y

so the normal to the velocity vector is orthogonal to the streamline. This means that the velocity is tangent to the streamline and justifies the interpretation that the particle of fluid at
x y is moving in the direction of the streamline at that point. We therefore interpret streamlines as the trajectories of particles in the fluid. If we sat on a particular particle, we would ride along a streamline. For this reason, graphs of streamlines form a picture of the motions of the particles of fluid. The remainder of this section is devoted to some illustrations of these ideas.

1090

CHAPTER 25

Conformal Mappings

EXAMPLE 25.23

Let f
z = −Kei z, in which K is a positive constant and 0 ≤  ≤ 2. Write f
z = −K cos
 + i sin
 x + iy = −K x cos
 − y sin
 − iK y cos
 + x sin
 If f
z = 
x y + i
x y, then 
x y = −K x cos
 − y sin
 and 
x y = −K y cos
 + x sin
 Equipotential curves are graphs of 
x y = −K x cos
 − y sin
 = constant. Since K is constant, equipotential curves are graphs of x cos
 − y sin
 = k or y = cot
x + b in which b = k sec
 is constant. These are straight lines with slope cot(. Streamlines are graphs of 
x y = − tan
x + d straight lines with slope − tan
. These lines make an angle  −  with the positive real axis, as in Figure 25.41. These are the trajectories of the flow, which may be thought of as moving along these straight lines. The streamlines and equipotential curves are orthogonal, their slopes being negative reciprocals. Now compute f

z = −Kei = −Ke−i  This implies that the velocity is of constant magnitude K. In summary, f models a uniform flow with velocity of constant magnitude K and making an angle  −  with the positive real axis, as in Figure 25.41.

EXAMPLE 25.24

Consider the flow represented by the complex potential f
z = z2 . This function is differentiable for all z, but f

0 = 0 so the origin is a stagnation point. We will see what effect this has on the flow and determine the trajectories. With z = x + iy, f
z = x2 − y2 + 2ixy, so 
x y = x2 − y2

and 
x y = 2xy

Equipotential curves are hyperbolas x2 − y 2 = k

25.5 Complex Function Models of Plane Fluid Flow

1091

y

x 

FIGURE 25.42 Equipotential curves and streamlines of the flow with complex potential f
z = z2 .

FIGURE 25.41

Streamlines of the flow with complex potential f
z = −Kei z.

if k = 0. Streamlines are hyperbolas xy = c if c = 0. Some curves of these families are shown in Figure 25.42. If k = 0 the equipotential curves are graphs of x2 − y2 = 0, which are two straight lines y = x and y = −x through the origin. If c = 0 the streamlines are the axes x = 0 and y = 0. The velocity of the flow is f

z = 2z. f models a nonuniform flow having velocity of magnitude 2 z at z. We can envision this flow as a fluid moving along the streamlines. In any quadrant, the particles move along the hyperbolas xy = c, with the axes acting as barriers of the flow (think of sides of a container holding the fluid).

EXAMPLE 25.25

We will analyze the complex potential f
z =

iK Log
z 2

Here K is a positive number and Log
z denotes that branch of the logarithm defined by Log
z = 21 ln
x2 + y2  + i, where  is the argument of z lying in 0 ≤  < 2 for z = 0. If z = x + iy, then    iK 1  2 K iK 2 f
z = ln x + y + i = − + ln
x2 + y2  2 2 2 4 Now 
x y = −

K  2

and 
x y =

K ln
x2 + y2  4

Equipotential curves are graphs of  = constant, and these are half-lines from the origin making an angle  with the positive real axis. Streamlines are graphs of 
x y = constant, and these are circles about the origin. Since streamlines are trajectories of the fluid, particles are moving on circles about the origin. Some streamlines and equipotential curves are shown in Figure 25.43.

1092

CHAPTER 25

Conformal Mappings

Equipotential lines

y

y

x

x Streamlines FIGURE 25.43 Streamlines and equipotential lines of the flow having complex potential f
z =
iK/2Log
z.

FIGURE 25.44

Streamlines of the flow having complex potential f
z =
iK/2Log
z.

It is easy to check that f

z =
iK/2
1/z if z = 0. On a circle z = r, the magnitude of the velocity of the fluid is   K K 1 
  = f
z = 2 z 2r This velocity increases as r → 0, so we have particles of fluid swirling about the origin, with increasing velocity toward the center (origin) (Figure 25.44). The origin is a vortex of the flow. To calculate the circulation of the flow about the origin, write f

z = −

iK z y x iK 1 K iK =− = − = u + iv 2 z 2 z2 2 x2 + y2 2 x2 + y2

If ! is the circle of radius r about the origin, then on !, x = r cos
 and y = r sin
, so


2 K r sin
 K r cos
 udx + vdy =
−r sin
 −
r cos
 d 2 r 2 2 r 2 ! 0 K
2 = −d = −K 2 0 This is the value of the circulation on any circle about the origin. By a similar calculation, we get −vdx + udy = 0 !

so the origin is neither a source nor a sink. In this example we may restrict z > R and think of a solid cylinder about the origin as a barrier, with the fluid swirling about this cylinder (Figure 25.45).

EXAMPLE 25.26

We can interchange the roles of streamlines and equipotential curves in the preceding example by setting f
z = KLog
z with K a positive constant. Now f
z =

K ln
x2 + y2  + iK 2

25.5 Complex Function Models of Plane Fluid Flow

1093

y

y Streamlines

x

x

z  R

Equipotential curves

FIGURE 25.45 Flow around a cylindrical barrier of radius R.

FIGURE 25.46 Equipotential curves and streamlines for the potential f
z = K Log
z.

so K ln
x2 + y2  and 
x y = K 2 The equipotential curves are circles about the origin and the streamlines are half-lines emanating from the origin (Figure 25.46). As they must, these circles and lines form orthogonal families of curves. The velocity of this flow is 
x y =

f

z =

x K y =K 2 + iK 2 = u + iv z x + y2 x + y2

Let ! be a circle of radius r about the origin. Now we find that udx + vdy = 0 !

and

!

−vdx + udy = 2K

The origin is a source of strength 2K. We can think of particles of fluid streaming out from the origin, moving along straight lines with decreasing velocity as their distance from the origin increases.

EXAMPLE 25.27

We will model flow around an elliptical barrier. From Example 25.25, the complex potential f
z =
iK/2Log
z for z > R models flow with circulation −K about a cylindrical barrier of radius R about the origin. To model flow about an elliptical barrier, conformally map the circle z = R to an ellipse. For this, consider the mapping a2  z in which a is a positive constant. This is called a Joukowski transformation, and is used in analyzing fluid flow around airplane wings because of the different images of the circle that result by making different choices of a. Let z = x + iy and w = X + iY . We find that the circle x2 + y2 = R2 is mapped to the ellipse w = z+

Y2 X2 + = R2  1 +
a/R2 1 −
a/R2 provided that a = R. This ellipse is shown in Figure 25.47. If a = R, the circle maps to

−2a 2a on the real axis.

1094

CHAPTER 25

Conformal Mappings z

2

w  z  az

w

z  R FIGURE 25.47 Joukowski transformation mapping a circle to an ellipse.

Solve for z in the Joukowski transformation. As a quadratic equation, this yields two solutions, and we choose √ w + w2 − 4a2  z= 2 Compose this mapping with the complex potential function for the circular barrier in Example 25.25 We get   √ iK w + w2 − 4a2 F
w = f
z
w = Log  2 2 This is the complex potential for flow in the w plane about an elliptical barrier if R > a, and about the flat plate −2a ≤ X ≤ 2a, Y = 0 if R = a. We will conclude this section with an application of complex integration to fluid flow. Suppose f is a complex potential for a flow about a barrier whose boundary is a closed path !. Let the thrust of the fluid outside the barrier be the vector Ai + Bj. Then a theorem of Blasius asserts that 1 A − iB = i f

z2 dz 2 ! in which  is the constant density of the fluid. Further, the moment of the thrust about the origin is given by   1
2 Re −  z f
z dz  2 ! In practice, these integrals are usually evaluated using the residue theorem.

SECTION 25.5

PROBLEMS

1. Analyze the flow given by the complex potential f
z = az, in which a is a nonzero complex constant. Sketch some equipotential curves and streamlines, determine the velocity, and whether the flow has any sources or sinks. 2. Analyze the flow having complex potential f
z = z3 . Sketch some equipotential curves and streamlines. 3. Sketch some equipotential curves and streamlines for the flow having potential f
z = cos
z. 4. Sketch some equipotential curves and streamlines for the flow having potential f
z = z + iz2 .

5. Analyze the flow having potential f
z = KLog
z − z0 , in which K is a nonzero real constant and z0 is a given complex number. Show that z0 is a source for this flow if K > 0 and a sink if K < 0. Sketch some equipotential curves and streamlines of this flow. 6. Analyze  the flow having potential f
z = z−a KLog , where K is a nonzero real number z−b and a and b are distinct complex numbers. Sketch some equipotential curves and streamlines for this flow.

25.5 Complex Function Models of Plane Fluid Flow   1 7. Let f
z = k z + , with k a nonzero real constant. z Sketch some equipotential curves and streamlines for this flow. Show that f models flow around the upper half of the unit circle. 8. Let   m − ik z−a f
z = Log  2 z−b in which m and k are nonzero real numbers and a and b are distinct complex numbers. Show that this flow has a source or sink of strength m and a vortex of strength k at both a and b. (A point combining properties of a source (or sink) and a vortex is called a spiral vortex). Sketch some equipotential curves and streamlines for this flow. 9. Analyze the flow having potential   1 ib f
z = k z + + Log
z z 2

1095

in which k and b are nonzero real constants. Sketch some equipotential curves and streamlines for this flow. 10. Analyze the flow having potential  √  √ 2z − ia 3 f
z = iKa 3Log √  2z + ia 3 with K and a positive constants. Show that this potential models an irrotational flow around a cylinder 4x2 + 4
y − a2 = a2 with a flat boundary along the y axis. Sketch some equipotential curves and streamlines for this flow. 11. Use Blasius’s theorem to show that the force per unit width on√the cylinder in Problem 10 has vertical component 2 3aK 2 , with  the constant density of the fluid.

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PA RT

8 Probability and Statistics

CHAPTER 26 Counting and Probability

CHAPTER 27 Statistics

V

XV

I II

XIII

III

XII

XI

XI

IV V

IX VI

IX VIII

VII Few areas of expertise have as profound an influence on us as statistics. Manufactured products are tested for reliability by performing statistical analyses of random samples. Risk analysis is an application of statistics. Statistics determine costs and availability of insurance, whether a new drug will receive FDA approval for distribution, which television programs will be shown and when, boundaries of congressional districts in the United States, the distribution of some kinds of federal aid, public health programs, and on and on. Professionals in business, medicine, science, engineering, and other areas use statistical analyses because they provide rigorous tools for analyzing information and drawing conclusions with measurable degrees of confidence. Understanding and using ideas from statistics require some knowledge of probability, which in turn makes use of counting techniques. Based on this chain of dependence, we will begin this part with counting, followed by probability, and then statistics.

1097

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CHAPTER

26

COMBINATIONS, PERMUTATIONS, MULTIPLICATION PRINCIPLE, PROBABILITIES, CONDITIONAL PROBABILITIES, BAYES’S THEOREM, TREE DIAGRAMS

Counting and Probability

Counting may sound simple enough. However, counting problems can be subtle. For example, how many lowercase, unordered nine-letter codes can be formed from the English alphabet if each vowel can be used no more than once? We will begin with techniques for solving the kinds of counting problems that we will encounter in statistics.

26.1

The Multiplication Principle Suppose we have some process that proceeds in n independent stages. Independent means that the outcome of one stage is not influenced by the outcome of the others. How many different ways can the entire process be carried out? As a simple example, suppose we are designing a car (the process), and we have four door designs and six fender designs. Door and fender selection are the two stages. How many different cars can we design from these choices, if each door looks good with each fender, and two cars are different if they differ in either a fender, a door, or both? This problem is easily solved. With the first door, we can use any of six fenders for six possibilities. With the second door, we can choose any of six fenders for six more cars. Similarly, there are six fenders that can go with the third door and six with the fourth. The total number of car possibilities is 6 + 6 + 6 + 6 = 24 Or, since there are six fenders with each of four doors, the number of cars is 4 · 6 = 24 To take a slightly more complicated example, suppose we have four doors, six fender styles, and nine hood shapes. Now how many cars can be formed? Since we can form 24 cars from the doors and fenders, and each can go with one of the hoods, there are 24 · 9 = 216 cars that can be made in this way. The answer in this case is 4 · 6 · 9. 1099

1100

CHAPTER 26

Counting and Probability

These examples suggest a general principle. The Multiplication Principle Suppose a process consists of n independent stages, or steps. Suppose stage j can be done in sj ways. Then the total number of ways of carrying out the process is the product s1 s2    sn  One way to look at this result is to envision n boxes, in each of which any of a certain number of objects (or choices) can be put. s1 s2 s3    sn−1 sn  The total number of ways of picking one object from each box is the product of the number of objects in each box.

EXAMPLE 26.1

A person flips a coin nine times. How many outcomes are there? We want to count something, namely the number of outcomes. However, before we count something, we should be clear on what we are counting. What is an outcome? An outcome here is the result of nine flips, hence will be a string of nine symbols, with the jth symbol an H (if toss j came up a head) or a T (if toss j was a tail). A typical outcome might be T H H T H H H T H The number of outcomes is the number of different nine letter strings, with each letter a T or an H. Think of nine empty boxes. There are two possible results we can put in the first box, namely T or H. There are also two for the second, and so on. Further, these flips are independent. The result of any flip does not influence the result of any other flip. The total number of possible outcomes of nine flips is therefore two multiplied by itself nine times, or 29 , which equals 512. This is an application of the multiplication principle, with n = 9 and each sj = 2. Similarly, the number of possible outcomes with thirty flips is 230 = 1 073 741 824

EXAMPLE 26.2

A game consists of flipping a coin four times, then rolling five dice. What is the number of outcomes? Here the activity or experiment has nine stages. The first four are all coin tosses, each with two possible outcomes. The last five are dice tosses, each with six outcomes. The number of possible outcomes is 2 · 2 · 2 · 2 · 6 · 6 · 6 · 6 · 6, or 24 65 , which equals 124 416. In the formalism of the multiplication principle, n = 9 and s1 = s2 = s3 = s4 = 2, while s5 = s6 = s7 = s8 = s9 = 6.

EXAMPLE 26.3

We want to form identification codes by choosing seven integers from 1 through 9, inclusive. How many ways can this be done?

26.1 The Multiplication Principle

1101

This problem has a wrinkle to it. If we choose the seven integers with replacement (that is, each can be any integer from 1 through 9, inclusively), then the total number of outcomes is 97 or 4 782 969. If, however, the seven integers are chosen without replacement, then every time we pick one it is used up, and we cannot choose it again. This means that there are nine choices for the first integer picked, but only eight for the second, seven for the third, and so on. Now the number of possible outcomes is 9 · 8 · 7 · 6 · 5 · 4 · 3 or 181 440.

SECTION 26.1

PROBLEMS

1. A game consists of flipping a coin seven times and rolling a die five times. (a) What is a typical outcome? (b) How many possible outcomes are there of this game? 2. A room has forty chests of drawers, and each chest has six drawers. A worker is assigned the task of inspecting each drawer in each chest. If each chest and each drawer take one second, how many seconds will it take the worker to complete this task? 3. A game consists of rolling ten dice, then choosing an integer falling between 7 and 21, inclusively. How many different outcomes are there to this game? 4. A man’s outfit for the day consists of a pair of shoes, a pair of trousers, a shirt, tie, jacket, sweater (but not both), and coat and a hat. He has seven pairs of shoes, twelve pairs of trousers, fifteen shirts, ten ties, three jackets, four sweaters, two coats and four hats. How many different outfits are possible, assuming that all the pairings making up an outfit are compatible (so, for example, any shirt looks good with any sweater)? How does the answer change, if at all, if an outfit consists of a pair of shoes, a pair of trousers, a shirt, tie, jacket and/or sweater, and a coat and/or hat? 5. ID codes for a company’s employees are to be formed as strings of six symbols. Each of the first three symbols can be any of the integers 1, 3 , 7, or 9, and each of the last three can be any of the letters A, C , D,

E or K. The company has 3 300 employees. Determine whether it is possible to assign each employee a different code in this way. 6. A factory produces tops for coffee tables. For each top, any of six colored panels are selected, then any of fifteen designer patterns is stamped on the selected panel. Finally, the underside of the panel is given any of four different finishes. How many different coffee table tops can be produced in this way? 7. A woman can fill out a set of golf clubs by choosing any one of seven putters, any one of fifteen drivers, any one of eight wedges, and any one of seventeen irons. How many different sets are there? 8. A carnival game consists of a large wheel with the numbers 1 through 50 spaced equally around the perimeter. A pointer is pivoted at the center and is spun by the operator, eventually coming to land on one of the numbers. The counter is spun four times. What does an outcome of these four spins look like, and how many different possible outcomes are there? 9. A committee of four is to be formed on the city council. The first member must be chosen from among seven residents of south side, the second from among eleven residents of city center, the third from among six residents of the park area, and the fourth from among fourteen residents of north beach. How many different committes can be formed? 10. Six coins are each flipped twenty times. What does an outcome look like, and how many different outcomes are there?

CHAPTER 26

1102

26.2

Counting and Probability

Permutations Suppose we have a collection of objects, and we want to arrange or list them in some order. How many ways are there to do this? This question has an important but simple answer. Think of having n boxes. We can put any of the numbers 1 2     n object in the first box, so there are n ways to fill this box. With one object chosen, there are n − 1 ways to fill the second box, then n − 2 ways to pick an object for the third box, and so on. Continue in this way. For the third from the last box, there are three remaining objects to choose from. For the next to last, there are two to choose from. Finally, there is only one object left, and that goes into the last box. By the multiplication principle, the total number of ways of making these choices is n ·
n − 1 ·
n − 2    3 · 2 · 1 which is the product of the integers from 1 through n. Our conclusion is: The number of ways of ordering n objects is the product of the integers from 1 through n, inclusive.

An ordering of objects is called a permutation of these objects. Thus, the number of permutations of three objects is 6, the number of permutations of four objects is 24, and the number of permutations of n objects is 1 · 2 · 3    n. This product occurs frequently in counting and other contexts (for example, coefficients in Taylor expansions), and is denoted n! (read “n factorial”). In this language: The number of permutations on n objects is n!

Here is a short list of factorials: n

n!

1 2 3 4 5 6 10 15 20 30 50

1 2 6 24 120 720 3,628,800 13077 × 1012 24329 × 1018 26525 × 1032 30414 × 1064

Factorials grow at an extremely rapid rate, 10! is already well over three million.

EXAMPLE 26.4

A lottery entry involves choosing any ordering of the objects a, b, c, d, e, and then any ordering of the objects A, B, C, D. How many possible outcomes are there in this lottery? Think of the game as consisting of two stages, the first choosing any ordering of a, b, c, d, e, the second choosing an ordering of A, B, C, D. There are 5! = 120 ways of doing the first step, and 4! = 24 ways of doing the second step. Further, the two steps are independent.

26.2 Permutations

1103

By the multiplication principle, there are 120 · 24 = 2880 possible outcomes. For example, one outcome is b c a d e B A D C Suppose we require that a be in the first position, and A in the last. How many outcomes are possible with this requirement? Think of nine boxes, with a already placed in the first and A in the last: a    A In the second through fifth places, there can be any arrangement of b, c, d and e. There are 4! such arrangements. In the sixth through eighth places, there can be any arrangement of B, C and D. There are 3! = 6 such arrangements. There are therefore 4!3! =
24
6 = 144 outcomes.

EXAMPLE 26.5

Suppose we are playing a poker game in which each hand has seven cards, but the order in which the cards is played is important. A particular player is dealt a jack, a king, a four, a two, a ten, a seven and a queen. How many different ordered hands can the player form? Since there are seven different cards, there are 7! = 5040 different orderings. This is the number of different hands that can be formed. The actual value of the cards is irrelevant information in reaching this conclusion. Given any seven different cards, there are 5040 orderings of them.

SECTION 26.2

PROBLEMS

1. How many ways can the first nine letters of the alphabet be arranged in different orders?

5. (a) How many arrangements are there of the objects a b c d f g h?

2. How many different codes can be formed from the lower case letters of the English alphabet, if a code consists of seventeen distinct letters, with different orders counting as different codes?

(b) How many arrangements are there if we insist on using only those beginning with a?

3. An ID number for each employee in a certain company consists of a string of nine numbers, consisting of an arrangements of the integers 1 through 9. How many different ID numbers can be formed?

6. We want to form ID numbers by using n distinct symbols and allowing any order for their arrangement. How large must n be to accommodate 20 000 people? How many for 1 000 000 people?

4. Suppose we have five symbols available, and these are %, $, #, ∗ and &. One plan is to form passwords by using different orderings of these five symbols. A second plan is to use ordered strings chosen from these five symbols, in which any symbol can be used one or more times. How many different passwords are possible under each plan?

7. Letters a through f are to be arranged in a three by two pattern, with one symbol in each location. For example,

(c) How many arrangements are there if a must be first and g fifth in the list?

da bf ce

CHAPTER 26

1104

Counting and Probability

is one such pattern. How many patterns can be formed in this way?

11. How many ways are there of choosing, with order, the even integers out of the integers 1 through 12?

8. Bowling pins are traditionally arranged in a triangle of fifteen identical pins. Suppose we decide to number the pins 1 through 15. How many different arrangements are there of these pins?

12. We have seen that n! increases at a fast pace as n is chosen larger. Because factorials arise in many contexts, there are formulas that can be used to approximate n! when n is large, but which can be computed in fewer steps than it takes to compute n! as the product of the integers from 1 through n. Stirling’s formula states that  n n √ n! ≈ 2n e with accuracy improving as n is chosen larger. To test this supposition, make a three column table, having values of n for n = 1 2    20 in the first column. In the second column, compute n!, and in the third column, the Stirling approximation. Compute the percentage error in the approximation in each case.

9. A lottery is run as follows. Twelve slips of paper are placed in a bowl. Each slip has a different symbol on it. A player makes an ordered list of these twelve symbols, and wins if they are drawn from the bowl in this order. How many possible different outcomes are there of this lottery? 10. The letters a through l are to be arranged in some order. How many arrangements are there that have a in the second place, d in the fifth place and k in the seventh place?

26.3

Choosing r Objects from n Objects Suppose we have n objects, and we want to choose r of them, with 1 ≤ r ≤ n. How many ways are there to do this? The answer depends on whether or not we take order into account. First consider the case that we do.

26.3.1

r Objects from n Objects, with Order

Given n objects, how many ways are there to make an ordered list of r of these objects? Using the multiplication principle, it is not difficult to derive a general formula for this number.

THEOREM 26.1

Let r and n be positive integers, with 1 ≤ r ≤ n. Then, the number of ways of picking r objects from n objects, taking order into account, is n! 
n − r! To understand this conclusion, envision r boxes. There are n choices of objects to put in box 1, leaving n − 1 choices for box 2, then n − 2 for box 3, and so on, until finally there are n − r + 1 objects left from which to pick one for box r. By the multiplication principle, the number of ways of making these choices of r objects from the given n objects is the product n
n − 1
n − 2   
n − r + 1

26.3 Choosing r Objects from n Objects

1105

This product is
n − r + 1    n = =

1 · 2 · · · ·
n − r − 1
n − r
n − r + 1 · · · · · n 1 · 2 · · ·
n − r n! 
n − r!

This number is often denoted n Pr , which is an abbreviation for “the number of permutations of r objects chosen from n objects.” Thus, n Pr

=

n! 
n − r!

For example, if n = 4 and r = 2, then =

4 P2

4! = 4 · 3 = 12 2!

This is the number of ways of choosing two ordered objects from four objects. As a convenience, define 0! = 1. In the present context, this makes sense. Suppose, for example, we want to compute 6 P6 . This is the number of ways of choosing six ordered objects from six objects. But this is exactly the number of permutations of six objects, which we know equals 6! Now compute 6 P6

=

6! 6! = = 6!
6 − 6! 0!

EXAMPLE 26.6

An election is being held and there are eight nominations for three offices. A ballot consists of three lines. The name filled in first receives one vote for president, the second name receives one vote for vice president, and the third name receives one vote for janitor. No name can be duplicated on any ballot. How many different ballots are possible in this election? Obviously the order in which candidate names are listed is significant here. The number of ballots is the number of ways of choosing three from eight, with order: 8 P3

=

8! = 6 · 7 · 8 = 336 5!

EXAMPLE 26.7

An ID code is to consist of eight different lowercase letters of the English alphabet, written in some order, together with an ordered string of any five distinct integers, chosen from 1 through 9, inclusively. How many different codes are there? There are 26 P8 ways of choosing eight of the 26 letters of the alphabet, with order, and 9 P5 ways of choosing five of the integers 1 through 9, with order. By the multiplication principle, the number of codes is 26! 9!
26 P8 
9 P5  = 18! 4! = 19 · 20    26 · 5 · 6    9 = 95242
1014 

1106

CHAPTER 26

26.3.2

Counting and Probability

r Objects from n Objects, without Order

Sometimes we want the number of ways of choosing r objects out of n objects, but the order in which the selected objects are listed does not matter. This number is denoted n Cr , and is called the number of combinations of r objects chosen from n objects. Whenever the word “ combination” is used, order is not taken into account. When “ permutation” is used, order is a factor. We already know how to compute n Cr . Given n and r, n Pr differs from n Cr only in taking the orderings into account. There are r! ways of ordering any r of the n objects. Therefore each combination of r of the n objects counted in n Cr gives rise to r! different orderings of these objects, and these are all included in the number n Pr . This means that n Pr

= r!n Cr 

Therefore, n Cr

=

1 n!  n Pr = r! r!
n − r!

For example, suppose we want to choose two objects out of four. If order is taken into account, then there are twelve possibilities, since 4 P2

=

4! = 12 2!

Explicitly, the possible choices are ab ac ad bc bd cd ba ca da cb db dc If order is not to be taken into account, then ab and ba are counted as one choice. Similarly, ac and ca are the same, and so on. There are six distinct choices (without order) of two of these four letters, and they are ab ac ad bc bd cd Without having this explicit list, we could still know that there are six such combinations by computing 4 C2

=

3·4 4! = = 6 2!2! 2

EXAMPLE 26.8

Elections are being held to form a committee to consist of four club members. There are seven nominations. A ballot consists of a list of four of the nominees, without regard to order, because the committee has no hierarchy. How many different ballots are there? Discount ordering and look for combinations, not permutations, of four members chosen from the seven nominees. The number of these combinations is 2·3·4·5·6·7 7! = = 35 7 C4 = 4!3! 2·3·4·2·3

EXAMPLE 26.9

Consider a poker game in which seven cards are dealt face down. Bets are placed after all hands have been dealt, at which time each player has seven cards making up the hand. What is the total number of possible hands?

26.3 Choosing r Objects from n Objects

1107

Assume a standard 52 card deck. Then the number of hands is the number of ways of choosing seven from 52 cards, without regard to order: 52 C7

=

52! 52! = 7!
52 − 7! 7!45!

=

46 · 47 · 48 · 49 · 50 · 51 · 52 = 133 784 560 2·3·4·5·6·7

Order does not matter here, because bets are made only after each player receives all seven cards, and nothing would change if the same seven cards were dealt in a different order. Now suppose each player is dealt one card face up, and then each player is dealt one card face down, at which time bets are made. Then each player is dealt one more card face down, followed by another round of betting. This continues until each player has seven cards (so six are face down, and there are six rounds of betting). Notice the difference between this and the first game. In the first game, the player only cares about which seven cards have been dealt, and the order made no difference. In the second game, order matters. If a queen is dealt first to a player (for all to see) and a second queen is dealt as the second card (for only the player to see), the player might bet in the early rounds differently than if the second queen is dealt in the last round. And opponents might bet differently if both queens had been dealt face down. Now order is vital to betting strategy. With order, the number of possible hands is 52 P7

=

52! = 46 · 47 · 48 · 49 · 50 · 51 · 52 = 674 274 182 400 45!

It is routine to check that this is 7! times Often n Cr is written

52 C7 .

n r

and

. For example,   6 6! = = 15 2!4! 2   5 5! = = 10 2 2!3!

This notation is often seen in connection with binomial expansions. If x and y are numbers, and n is a positive integer, then
x + yn =

n   n k=0

k

xk yn−k 

  Because of their appearance as coefficients in a binomial expansion, the numbers nk are called binomial coefficients. The coefficient of xk (which is the same as the coefficient of xn−k ) is the number of ways of choosing k objects from n objects, disregarding the order of the choice. This makes sense, because if we choose k from n objects, then we have automatically also chosen n − k of the n objects.

26.3.3

Tree Diagrams

In working with permutations and combinations, it is often handy to use a systematic device called a tree diagram. An illustration will clarify what this is.

1108

CHAPTER 26

Counting and Probability

Consider strings of symbols of length 5, using either an H or a T in each place. For example, HHTHT

and HTTTT

are two such strings. There are 25 = 32 strings. Now, suppose we want the number of strings having exactly two T ’s (and therefore three H’s). Think of a string of five boxes. We want to pick out two to put T in, the rest automatically getting an H. There are 5 C3 = 10 ways of doing this. Order is irrelevant here because, once we pick two boxes, we put an identical T in each, then an identical H in the other three. None of this involves a new idea. But now suppose we want to actually list all ten of these strings. Here is a systematic way to do this. Put a dot on a piece of paper to represent a starting point. Next, place two dots to the right (or left, if you prefer), one for H and the other for T , with a line (or edge, or branch) from the starting dot to each of these. So far the diagram displays the outcomes of the first flip. Next draw two lines, ending in dots for H and T , from each of the first H and T dots. The diagram so far, Figure 26.1, shows outcomes of the first two flips. Continue, putting two lines from each of the four end dots drawn so far, with each line ending with an H or a T dot (Figure 26.2). The pattern is obvious, and we carry out two more steps (two lines from each H or T , then repeat once more), until strings of five edges form along each path from the starting dot to the the right-most dots (Figure 26.3). The resulting configuration is called a tree diagram, and it is useful for reading information about

H H T H

H

T

T

H

H H

H

T

T T

H H T

T T

T

FIGURE 26.1

FIGURE 26.2

Outcomes of the first two coin flips.

three coin flips.

Outcomes of the first

26.3 Choosing r Objects from n Objects H T

H

H T

T

H

1109

H

H

H

T

T H H

T H

H T H T

T

H T H

T T

H T

T H T H H

T H

T

T

H H

H T H T H

H T T H

T H

T H T

T H

T T

H T H

T T FIGURE 26.3 Outcomes of five coin flips.

the strings. Each path of five edges, from the starting dot to an end dot on the right, details one complete sequence of five flips of the coin. By following various paths, we can, for example, easily pick out exactly those having two T ’s and three H’s: HHHTT HHTHT HHTTH HTHHT HTHTH HTTHH THHHT THHTH THTHH TTHHH Although we knew without the tree diagram that there would be ten of these outcomes, the tree diagram is a systematic way of listing them all, if that is needed.

1110

CHAPTER 26

Counting and Probability 1 23 4

1 2

3 4 5

H 6

H 6 5

43

2

1 6

4 5

5 6 1 2 3 4 5 6 1 2 3 4 5 6 1 2 3 4 5 6 1 2 3

1 2 3 4

5

6

1 2

3 4

1

5 6

2

1 2 3 4 5 6

3

4

1 2 3 4

5 1 6 2 3

T 6

5

4 6 6 5 4

3

2

5

1 1

1

2 3 4 5 6

1

2

H

2 3 4 5 6

3

1 2

T

4 3

1

1 2

T

6

5

4 5 6

3

23 4 5 6 1 2

6

5

1 2 3 4

6

4 5 6

65

FIGURE 26.4

twice).

4

5

1 2 3 4

6 3 4 5 6 6

1 2 3

1 2 3 4 5

6

5

1 2 3 4

5

4

1 2 3

1 2 3 4 5 6

Outcomes of flipping two coins or one coin twice), then rolling two dice (or one die

26.3 Choosing r Objects from n Objects

1111

EXAMPLE 26.10

A game consists of flipping two coins, then rolling two dice. An outcome will consist of a string of length four, the first two symbols recording the results of the coin flips (H or T ) and the last two the numbers that come up on the dice
1 through 6). There are 2 · 2 · 6 · 6 (or 72) outcomes. Figure 26.4 shows a tree diagram for this game. By following all of the paths from the left starting dot to the end dots on the right, we can read all of the outcomes explicitly. For example, all of the outcomes in which the two dice come up both 3, or both 4, or one 3 and one 4, can be easily read from the diagram: HH33 HT 33 TH33 TT 33 HH44 HT 44 TH44 TT 44 HH34 HT 34 TH34 TT 34 HH43 HT 43 TH43 TT 43

SECTION 26.3

PROBLEMS

1. Seven members of an audience of 25 people are to be chosen to win a prize. The first name drawn will win half of the planet, the second name drawn will win an airplane, the third name drawn a new house, and so on down to the last name drawn, which will win fifty cents worth of merchandise at the nearest drug store. How many different possible outcomes are there of this drawing? 2. There are five positions open on the board of a swimming club, and sixteen people are eligible for election to the board. A ballot consists of a list of names of five of the eligible members, with order being important because new board members are assigned positions of decreasing importance, depending on how far down the list a person is selected. How many different ballots are possible in this election? 3. A car dealer has a lottery. First, 22 names are selected at random from a data base of past customers. Six of these names will be winners. The first name chosen can go through the lot and pick any car. The second person can do the same, but perhaps the best car has already been chosen by the first person. The third person also goes through the lot, taking third pick, and so on. How many different possibilities are there for lists of people to go through the lot and make choices? 4. A game consists of choosing and listing, in order, three numbers from the integers 1 2     20.

(a) How many different choices are there? (b) What percentage of these begins with the number 4? Would this answer change if the first number is 17 instead of 4? (c) What percentage of these choices ends with the number 9? Would this number be the same if the last number were 11 instead of 9? (d) How many choices are there if the first number must be 3 and the last number 15? 5. How many different ways can a ten card hand be dealt from a standard 52 card deck, if the order in which the cards are dealt is unimportant? What is the number if the order of the deal is significant? 6. How many different nine man lineups can be formed from a roster of 17 players, if order of selection is not a factor? How many can be formed if order makes a difference? 7. How many different ways can four drumsticks be chosen from a barrel containing twenty drumsticks, if order of selection does not matter? 8. A company is selecting twelve of its forty employees to lay off. How many different ways can such a selection be made, if the order of choice is unimportant? 9. A carnival game begins with the contestant drawing five cards from a standard deck. Contrast the number

CHAPTER 26

1112

Counting and Probability

of different ways this can be done if order is important, with the number of different hands that can be drawn if order does not count. 10. How many different ways can a seven person committee be selected from a group of seventeen people, (a) if the order of selection is important, and (b) if the order does not matter? 11. Let n and k be integers with 0 ≤ k ≤ n. Show that n  n =  k n−k 12. Five coins are tossed. (a) Draw a tree diagram for this experiment.

13. Three dice are rolled. (a) Draw a tree diagram for this experiment. (b) List all outcomes in which the sum of the dice is even. (c) List all outcomes in which the first and third die are odd. (d) List all outcomes in which the first die comes up 2 or 5, and the last die is 3, 4 or 6. 14. Two coins are tossed and then and two dice are rolled. (a) Draw a tree diagram for this experiment.

(b) List all outcomes in which the second and third coins are heads.

(b) List all outcomes having a head on the first coin and the sum of the dice is odd.

(c) List all outcomes in which the first coin comes up heads and the fourth, tails.

(c) List all outcomes in which the first die is a 1 or 5.

(d) List all outcomes in which two coins come up heads and two come up tails.

(d) List all outcomes in which both coins come up tails and the second die is 1, 3 or 4.

26.4

Events and Sample Spaces In this section we begin to establish a framework within which we will be able to compute probabilities of certain kinds of events. In order to isolate the essential issues, begin with a simple probability problem that we can all solve. Flip an honest coin. What is the probability that it comes up heads? Everyone knows that this probability is 1/2. Although this probability is sometimes stated as 50%, we prefer to specify any probability as a number from 0 to 1 inclusive. Why is this obvious? Perhaps without explicit verbalization, 1/2 seems natural for this probability because a head is one of two equally likely outcomes, heads or tails. Formally, Probability (head) =

number of ways the coin can come up heads 1 =  total number of possible outcomes of one flip 2

Try this reasoning on a similar but slightly more complicated problem. Roll a single die. What is the probability that it comes up 3? It is common to guess that this probability must be 1/6. Why? There are six possible outcomes (six faces on the die), and they are all equally likely (honest die), so 3 will likely come up one in six times, for a probability of 1/6. Again, Probability (rolling 3) =

number of ways 3 can come up 1 =  total number of possible outcomes of one roll 6

Notice a common thread in both of these probabilities. In each, we performed an experiment (flip a coin, roll a die), and there were certain outcomes. The probability of an outcome was the number of ways this outcome could occur, divided by the total number of possible outcomes. These simple examples are useful in suggesting an approach that will enable us to handle more complicated probability questions.

26.4 Events and Sample Spaces

1113

The fundamental setting for any question in probability is an experiment, by which we mean simply something that is done. It could be a coin flip, a roll of a die, pulling an item from a production line, picking twelve cards from a deck, choosing a marble from a jar, guessing a lottery number, choosing a patient in a clinical trial, or almost any other action that has a finite number of outcomes. We will not consider experiments with infinitely many outcomes. We assume that each experiment we deal with has known outcomes. The number of outcomes may be very large, so we do not want to have to actually make a list to compute a probability. But we do need to know what the outcomes are and at least be able to describe them all by some rule or rules. We will also have to be able to count them. For example, suppose an experiment is to choose an integer at least as large as 1 and no larger than ten hundred billion. This experiment has ten hundred billion possible outcomes, namely all of the integers in the specified range. It would take us a long time to actually write these all down in a list, and there is no need to do this. The important thing is that we know all the outcomes, what they look like, and how many there are.

DEFINITION 26.1

Sample Space

The set of all outcomes of an experiment is called its sample space.

EXAMPLE 26.11

An experiment is to flip two coins. Now there are four possible outcomes. If the sample space is called Q, then Q = HH TT HT TH These are all the outcomes. Notice that TH and HT are separate outcomes. There are two ways we can get a head and a tail—a head on the first coin and tail on the second or a head on the second coin and tail on the first.

EXAMPLE 26.12

Roll two dice. How may outcomes are there? We can record the result of rolling two dice as a pair
a b, where a is the number that comes up on the first die, and b is the number on the second die. There are six choices for a, and six for b, so by the multiplication principle, there are 6 · 6 = 36 outcomes. If we list them explicitly, the sample space W is W = 
1 1
1 2
1 3
1 4
1 5
1 6
2 1
2 2
2 3
2 4
2 5
2 6
3 1
3 2
3 3
3 4
3 5
3 6
4 1
4 2
4 3
4 4
4 5
4 6
5 1
5 2
5 3
5 4
5 5
5 6
6 1
6 2
6 3
6 4
6 5
6 6

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EXAMPLE 26.13

An experiment consists of choosing (without replacement) five different letters of the English alphabet and writing them in some order. What are the outcomes, and how many outcomes are there? Suppose we use only lowercase letters. Then there are 26 letters from which to choose five. Order counts because of the way the experiment was defined. Thus there are 26 P5

=

26! = 22 · 23 · 24 · 25 · 26 = 7893600 21!

outcomes, each of which is a string of different letters (because the letters were chosen without replacement). Typical outcomes are acdfg and gdfac. We know what an outcome looks like, and how many there are. Once we understand the concepts of an experiment and the resulting sample space, we can define an event for an experiment.

DEFINITION 26.2

Event

An event for an experiment is a set of outcomes.

The rationale for this definition is that we sometimes want the probability of something that is more complicated than a single outcome, and consists of several outcomes specified in some way.

EXAMPLE 26.14

Flip three coins. A natural way to record an outcome of this experiment is as a 3-character string of H’s and T ’s. Since each flip has two possible outcomes, there are 23 = 8 outcomes. This is a small enough number that we can actually write out the sample space explicitly: S = HHH HHT HTH HTT TTT TTH THT THH An event in this experiment is any collection of outcomes. The entire sample space S is an event. Another is E1 = HHH TTT HTH Another event is E2 = HHH TTT This is the event that all three flips come up the same. Still another event is E3 = HTT THT TTH the event that there is exactly one head.

26.4 Events and Sample Spaces

1115

EXAMPLE 26.15

Consider the experiment of rolling two dice. The sample space W consists of 36 outcomes, each of the form
a b, where a and b can independently be any integer from 1 through 6, inclusively. Suppose we are interested in the event E that the sum of the numbers that come up on the dice is 5. E consists of all outcomes
a b with a + b = 5: E = 
1 4
2 3
3 2
4 1 These are all the ways the dice can total 5. Suppose we want the event B consisting of all the ways we can roll a total of 2. This is the event B = 
1 1 There is only one way to roll a total of 2. We have the intuition that it is less likely to roll a total of 2 than a total of 5, because there is only one way to roll 2, but four ways to roll 5. It will be useful to have an empty event. This event is denoted ∅, and consists of no outcomes. To illustrate, in the last example, the event that the dice total 13 is empty, containing no outcomes, since this can never happen. We are now ready to compute probabilities. In doing this, we emphasize a crucial point, the difference between an outcome and an event. An outcome is something that occurs when the experiment is performed. An event is a collection of one or more outcomes (or even perhaps no outcomes). Outcomes are the building blocks of events.

SECTION 26.4

PROBLEMS

In each of Problems 1 through 10, an experiment is described. Describe a typical outcome, and count the number of outcomes. Determine the number of outcomes in each of the events defined for each experiment. It is not necessary to list all of the outcomes of an event explicitly.

4. Pick five distinct lowercase letters from the English alphabet, taking order into account. Event I consists of all outcomes beginning with zk. Event W consists of all outcomes having third letter e. Event P consists of all outcomes with chjk as the last four letters.

1. Roll four dice. Event A consists of all outcomes in which the dice total 9. Event B consists of all outcomes in which each die comes up 1 or 2.

5. Flip three coins and then roll three dice. Event C consists of all outcomes with three heads and the dice totalling 15. Event D consists of all outcomes with at least two heads and the dice totalling at least 15. Event E consists of all outcomes with exactly one head and only a five or six showing on the each die.

2. Draw seven cards from a standard 52 card deck, with order of the draw having no importance. Event W consists of all hands having only cards numbered 2 or 3. Event K consists of all hands having only kings and/or aces. Event M consists of all hands with exactly one ace. 3. Pick four distinct lowercase letters from the English alphabet, without regard to order. Event A consists of the outcomes having an a and a z. Event W consists of the outcomes having only letters between a and g, inclusively.

6. Pick an integer from 0 through 10 inclusive, then three letters from the English alphabet. M is the event that the number chosen is even. K is the event that the number is 3, 5 or 7, and one of the letters is w. 7. Five dice are rolled. U is the event that the dice total at least 28. K is the event that the dice total no more than 7.

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Counting and Probability jacks, four kings and three aces were drawn (together with other cards).

8. Twenty three coins are flipped. R is the event that at least twenty heads come up. Q is the event that the last nineteen flips all came up heads. 9. Sixteen cards are drawn (without regard to order) from a standard 52-card deck. Y is the event that every card drawn was a face card or ace. M is the event that four

26.5

10. A ten sided polygon has its sides numbered 1     10. Four sides are selected at random. Y is the event that the sides total exactly 14. P is the event that sides 2, 3 and 7 were picked.

The Probability of an Event Suppose some experiment has taken place, and we focus on a certain event. We want to assign a number, between 0 and 1, to this event in such a way as to give a measure of the likelihood, or probability, that this event will occur if we perform this experiment. In this scheme, probability 0 means that the event cannot occur, probability 1 means that it is certain, and numbers strictly between 0 and 1 are assigned to events that are possible, but not certain. Further, these numbers must be assigned so that larger probabilities correspond in a reasonable way to “more likely” events. To calculate the probability of an event, the following must be in place. 1. There must be an experiment and a sample space S = o1  o2      oN  with finitely many outcomes. We call oj the j th outcome. 2. For each outcome oj , there must be given, or determined, a number Pr
oj , with 0 < Pr
oj  < 1. Pr
oj  is called the probability of oj . 3. The sum of the probabilities of all the outcomes must equal 1: Pr
o1  + Pr
o2  + · · · + Pr
oN  = 1 The function Pr, operating on the outcomes, is called a probability function, or probability measure, for this experiment. When we have a probability function, the probability of an event E is denoted Pr
E, and is defined to be the sum of the probabilities of the outcomes in E: Pr
E = sum of probabilities of outcomes in the event E. For example if E = o3  o7  o10 , then Pr
E = Pr
o3  + Pr
o7  + Pr
o10  in which each of the probabilities of outcomes on the right is assumed known. The probability of the empty event, ∅, is defined to be zero, and the probability of the entire sample space (which is an event), is one: Pr
∅ = 0 Pr
S = 1 The fact that Pr
S = 1 is a necessary consequence of Property 3 of the probability function. S is itself an event containing all of the outcomes, and the probability of this event is by definition the sum of the probilities of all the outcomes in S, which must equal one. Finding a probability function for an experiment is now the central issue in computing probabilities of events.

26.5 The Probability of an Event

1117

EXAMPLE 26.16

Consider a single coin flip. We do not need all of this machinery for such a simple experiment, but it is instructive to see how it works here. This experiment is to flip a coin once. There are two outcomes, which we call H and T , so the sample space is S = H T. Here o1 = H and o2 = T and the number of outcomes is N = 2. We need to assign a probability to each outcome. If the coin is honest, then we expect a head to have the same likelihood of coming up a tail, so we set 1 Pr
H = Pr
T =  2 This is consistent with Property 3, since the sum of the probabilities of all the outcomes is one: 1 1 + = 1 2 2 In this experiment, there are only four events, namely the empty event ∅, H consisting of just H, T consisting of just T , and H T, which is the entire sample space. Pr
H + Pr
T =

EXAMPLE 26.17

We will redo the last example with a twist. Suppose the coin is dishonest, and on any flip, a head is twice as likely as a tail. If o1 = H and o2 = T , then Pr
o1  = 2 Pr
o2  Now, to be a probability measure, we require that Pr
o1  + Pr
o2  = 1 These equations imply that 2 1 and Pr
T =  3 3 This is a legitimate probability function for this experiment, but it is different from that of the preceding example because the coin is dishonest in a known way. Pr
H =

We will now move to more complicated experiments, beginning with a class of special probability problems having equally likely outcomes. Suppose an experiment has N outcomes o1      oN , and each is equally likely: Pr
o1  = Pr
o2  = · · · = Pr
oN  Since Pr
o1  + Pr
o2  + · · · + Pr
oN  = N Pr
o1  = 1 then necessarily Pr
o1  =

1  N

Because all the outcomes are equally likely, then Pr
oj  =

1 N

for j = 1 2     N

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Counting and Probability

Such an experiment is called an equally likely outcome experiment. We also say that the sample space is an equally likely sample space. In such a case, if an event has k outcomes in it, then the probability of this event is the sum of the probabilities of the outcomes in the event, which is 1/N added to itself k times, or k/N . We will record these results for future use. THEOREM 26.2

In an equally likely outcome experiment with N outcomes, the probability of each outcome is 1/N . Further, if an event E has k outcomes, then Pr
E =

k  N

Another, sometimes useful, way of writing this result is that, in an equally likely outcome experiment, Probability of an Event E =

the number of outcomes in E the total number of outcomes of the experiment

=

the number of ways the event can occur  the total number of outcomes of the experiment

(26.1)

This reduces many probability calculations to counting problems.

EXAMPLE 26.18

Flip three coins. What is the probability of getting exactly two heads? There are 23 = 8 outcomes, each equally likely. The probability of any one outcome is 1/8. The event E of getting exactly two heads is E = 
H H T
H T H
T H H containing three outcomes. The probability of getting exactly three heads is therefore 3 Pr
E =  8 EXAMPLE 26.19

Three dice are rolled. What is the probability that the dice total 17? Each outcome is a triple
a b c, with a the number that came up on the first die, b that on the second, and c the number on the third die. Assuming honest dice, all outcomes are equally likely. There are N = 63 = 216 outcomes. The event E we are interested in (dice total 17) is E = 
6 6 5
6 5 6
5 6 6 There are k = 3 outcomes in E. By equation (26.1), Pr
E = 3/216

EXAMPLE 26.20

An experiment consists of flipping 25 honest coins. What is the probability that exactly seven coins come up heads?

26.5 The Probability of an Event

1119

This is an equally likely outcome experiment. We need the number of outcomes, and the number of outcomes in which there are exactly seven heads. Since there are 25 flips, each with two possible outcomes, the experiment has 225 = 33 554 432 outcomes. The event E we are interested in is that exactly seven tosses yield H (so the other 18 are all tails). The number of outcomes having exactly seven H’s is the number of ways of choosing seven out of 25 without regard to order, and this number is 25 C7

=

25! 25! = = 480 700 7!
25 − 7! 7!18!

Now we have both the numerator and denominator in equation (26.1), and Pr
E =

480700 ≈ 00143 33554432

What is the probability of coming up with at least 22 heads? Let F be this event. The key phrase here is “at least.” We obtain at least 22 heads if we roll exactly 22 heads, or exactly 23, or exactly 24, or all 25 heads. These are all different outcomes comprising parts of event F . We will count all of these outcomes: Number of outcomes with at least 22 heads = number of outcomes with exactly 22 heads + number of outcomes with exactly 23 heads + number of outcomes with exactly 24 heads + number of outcomes with exactly 25 heads =

25 C22 + 25 C23 + 25 C24 + 25 C25

=

25! 25! 25! 25! + + + = 2626 3!22! 2!23! 1!24! 25!

From equation (26.1), Pr
F =

2626 ≈ 0000078 33554432

As intuition might suggest, this event is not one to bet on.

EXAMPLE 26.21

Suppose a person is dealt five cards, without regard to order, from a standard deck. What is the probability that the hand has four aces in it? The total number of outcomes is the number of ways of choosing five cards from 52, 52! = 2598960 disregarding order. This number is 52 C5 = 5!47! This is the denominator in equation (26.1). Now we need the numerator. How many unordered five-card poker hands have exactly four aces? If four of the cards are aces, the fifth card can be any of the remaining 48 cards. Therefore the probability of being dealt four aces is Pr
four aces =

48  2598960

which is approximately 0000018. This is why, in westerns, a gambler showing four aces usually gets shot.

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We will consider experiments in which not all outcomes have the same probability when we have more machinery to deal with them efficiently.

SECTION 26.5

PROBLEMS

1. Suppose five honest coins are flipped. (a) Find the probability of getting exactly two heads. (b) Find the probability of getting at least two heads. 2. Roll four dice.

(a) What is the probability that all seven dice came up 3? (b) What is the probability that five dice came up 1, and two came up 4? (c) What is the probability that the sum of the numbers that came up is 26?

(a) What is the probability that exactly two 4’s come up?

(d) What is the probability that the sum of the numbers that came up is at least 26?

(b) What is the probability that exactly three 4’s come up?

7. Twenty balls in an urn are numbered 1 through 20. A blindfolded contestant draws five balls from the urn, with the order of the draw recorded.

(c) What is the probability that at least two 4’s come up? (d) What is the probability that the dice total 22?

(a) What is the probability that balls 1 2 3 4, and 5 were drawn (in this order)?

3. Two cards are selected from a standard deck. The order of the draw is unimportant.

(b) What is the probability that the number 3 ball was selected?

(a) What is the probability that both cards were kings?

(c) What is the probability that an even-numbered ball was drawn?

(b) What is the probability that neither card was an ace or face card? 4. Four letters are selected from the English alphabet (only lowercase is used). The order of the selection is recorded. (a) What is the probability that the first letter was q? (b) What is the probability that a and b are two of the letters? (c) What is the probability that the letters are a, b, d, z, in this order? 5. Eight bowling balls are in a bin. Two are defective, having been manufactured as cubes instead of the traditional spherical shape. A person uses a remote gripping device to pick three of the balls out of the bin, sight unseen. The order of the choice is unimportant. (a) What is the probability that none of the balls chosen was defective?

8. Seven drawers in a desk contain a 50-cent piece, while two other drawers contain a thousand dollar bill. A person is allowed to choose three drawers at random. The order of the draw is unimportant. (a) What is the probability that the person gets at least $1,000? (b) What is the probability that the person ends up with less than one dollar? (c) What is the probability that the payoff is $150? 9. Five cards are drawn, without regard to order, from a standard deck. (a) What is the probability that the hand contains at exactly one jack and exactly one king? (b) What is the probability that the hand contains at least two aces?

(b) What is the probability that exactly one defective ball was taken?

10. Twenty integers are chosen at random, and without regard to order, from the integers 0 to 100, inclusively.

(c) What is the probability that both defective balls were taken?

(a) What is the probability that all of the numbers chosen are larger than 79?

6. Seven pyramid-shaped (tetrahedron) dice are tossed. The faces on each pyramid are numbered 1 through 4.

(b) What is the probability that one of the numbers is 5?

26.6 Complementary Events

26.6

1121

Complementary Events If E is an event for some experiment, then E consists of certain outcomes. The outcomes that are not in E form the complement of E, denoted E C . This is another event. Notice that the complement of E C is E again:
E C C = E We may think of E C as the event “E does not occur”. The outcomes in E, together with those in E C , are all of the outcomes. Further, E and E C have no outcomes in common. If we listed all of the outcomes in E, and then all those in E C , we would list all the outcomes in the sample space S, with no repetitions. This means that Pr
S = Pr
E + Pr
E C  = 1 Another way of looking at this equation is that, for any event E, it is certain (probability 1) that either E occurs, or E does not occur (which is the same as E C occurs). If we write this equation in a slightly different way, we obtain the following. Principle of Complementarity: Pr
E = 1 − Pr
E C  This equation holds for E = S as well, since then E C = ∅, so Pr
E = Pr
S = 1 and Pr
E C  = Pr
∅ = 0. The principle of complementarity is useful because it offers us the choice of computing Pr
E or Pr
E C , and one might be easier than the other.

EXAMPLE 26.22

Roll three dice. There are 63 = 216 outcomes. Consider the event E that the dice come up with a total of at least 5. Compute Pr
E. One way to do this is by using equation (26.1): number of outcomes in E  216 Counting the number of outcomes in E is certainly possible, but it will be tedious. E contains all outcomes with the dice summing to 5 or more, meaning that the dice can total 5, 6, 7, 8,     18. We would have to count the number of ways each of these totals can occur. However, look at E C , which consists of all tosses of the three dice in which the total is 4 or less. There are many fewer such outcomes, making them easier to count. The only way the dice can total 4 is to come up 2 1 1 or 1 2 1 or 1 1 2 (three ways). The only way they can total less than 4 is to total 3, which can happen only one way (they all come up 1). Therefore E C has exactly four outcomes in it, so Pr
E =

Pr
E C  =

1 4 = 216 54

and by the principle of complementarity, 53 1 = ≈ 098 54 54 Rolling a total of at least 5 with three dice is not certain, but it is very likely. Pr
E = 1 −

CHAPTER 26

1122

SECTION 26.6

Counting and Probability

PROBLEMS

1. Seven dice are rolled. What is the probability that at least two of them come up 4? 2. Fourteen coins are tossed. What is the probability that at least three come up heads?

5. Four numbers are chosen at random and without regard to order from among the integers from 1 through 55, inclusively. What is the probability that at least one number is greater than 4?

3. Five cards are drawn from a standard deck, without regard to order. What is the probability that at least one is a face card, an ace, or numbered 4 or higher?

6. Six dice are rolled. What is the probability that they total at least 11?

4. Two coins are tossed and five dice are rolled. What is the probability that two heads and at least one 4 come up?

7. Here is the famous Birthday Problem. Suppose N people are in a room. What is the probability that at least two have the same birthday?

26.7

Conditional Probability Probability depends in subtle ways on how much is known. Here is an illustration of this. Suppose we pick a card at random from a deck of 52 cards. What is the probability that we chose a diamond? Since one fourth of the cards are diamonds, the probability of choosing a diamond is 1/4. However, suppose, as the card was drawn, it flashed for an instant in front of a mirror and we saw that the card’s suit is red. We immediately know that the card must be a diamond or heart. This narrows our choice down to one suit out two, not four, so the probability of a diamond is now 1/2. The additional knowledge that the suit was red has changed the sample space, which now consists of only those outcomes that are consistent with the additional information that the card chosen was red. Which is correct, 1/4 or 1/2? The answer is both are right, because these numbers are different probabilities. The first is a straightforward probability of picking a particular suit of card from the entire deck. The second is a conditional probability, in which we not only pick a card, but also have additional information that eliminates certain outcomes as possibilities, hence reduces the sample space. To put all of this on a firm footing that we can work with, suppose an experiment leads to outcomes and a sample space S. Single out a particular event U , having positive probability. (This will be the additional information of the conditional probability). If E is any event, define E ∩ U to be the event consisting of outcomes common to E and to U . For example, if E = a b c d e f g h and

U = k g h a w z

then E ∩ U = a g h This is also an event. E ∩ U is read “E intersect U ,” or “the intersection of E and U ,” or, in the probability context, “E and U .” Notice that E ∩ U = U ∩ E, since the act of taking all outcomes common to E and U does not depend on the order in which E and U are written. It is possible that there are no outcomes common to E and U , in which case E ∩ U = ∅, the empty event.

26.7 Conditional Probability

1123

Now imagine that the experiment is performed, but that, by some means (such as the mirror in the card experiment), we know that U occurs. The conditional probability of E, knowing U , is denoted Pr
E  U and is computed as Pr
E  U =

Pr
E ∩ U  Pr
U

(26.2)

By taking the probability of E ∩ U in the numerator, we are factoring in the information that U is assumed known, so we only look at outcomes in E that are also in U . In the denominator, we have the probability of U computed with respect to the original sample space of the experiment. In effect, we are considering only outcomes of E that are in U , and thinking of U as the new sample space. We will redo the card example to illustrate these ideas and how equation (26.2) works.

EXAMPLE 26.23

Pick a card out of a 52-card deck. What is the probability that a diamond is chosen? Since 13 of the 52 cards are diamonds, Pr
E =

13 1 =  52 4

Now let U be the event that the card is a diamond or heart. This is the information flashed in the mirror. In this example E ∩ U = E, because every event in E (card is a diamond) is already in U (card is a diamond or heart). Further, U consists of exactly half of the cards, so Pr
U = 1/2 (in the original experiment). Then Pr
E  U =

Pr
E ∩ U Pr
E 1/4 1 = = =  Pr
U Pr
U 1/2 2

consistent with intuition. Knowing that the card is a diamond or heart, the probability of picking a diamond is 1/2.

EXAMPLE 26.24

Toss three dice. What is the probability of rolling a total of 5? This is not a conditional probability, but we need this result for the conditional probability we will consider. The sample space consists of triples
a b c, in which each letter can be any integer from 1 through 6 inclusive. There are 63 = 216 outcomes. The outcomes in which the dice total 5 are
1 1 3
1 3 1
3 1 1
1 2 2
2 1 2
2 2 1 six in number. These constitute the event E: Roll a total of 5. Then Pr
E =

6 1 =  216 36

Now suppose a person sitting on the side knows that two dice are loaded, and must always come up 2. This means that at least two dice must come up 2 on any roll of all three dice. Let U consist of all outcomes
x 2 2
2 x 2

and
2 2 x

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in which x (for the honest die) can be any of 1, 2, 3, 4, 5, or 6. There are 16 such outcomes (note that when x = 2 all three triples are the same
2 2 2), so Pr
U =

2 16 =  216 27

Now E ∩ U = 1 2 2
2 1 2
2 2 1 the outcomes common to E and U . There are three such outcomes, so 3 1 =  216 72 The conditional probability that E will happen, knowing that U happens, is Pr
E ∩ U =

Pr
E  U =

Pr
E ∩ U 1/72 3 = =  Pr
U 2/27 16

This is the probability that the dice will sum to 5, as far as the person who knows about the two loaded dice is concerned. Notice that the probability of E (dice sum to 5) is only 1/36. However, to the person who knows that two dice always come up 2, the probability of the dice summing to 5 is considerably improved to 3/16. Knowing that two dice always come up 2 changes the sample space by eliminating certain outcomes. Examine equation (26.2) more closely. By using equation (26.1) twice, we can write Pr
E ∩ U =

number of outcomes common to E and U number of outcomes in S

and number of outcomes in U  number of outcomes in S From these equations a conditional probability can be computed as Pr
U =

Pr
E  U =

Pr
E ∩ U Pr
U

number of outcomes common to E and U number of outcomes in S = number of outcomes in U number of outcomes in S number of outcomes common to E and U  (26.3) = number of outcomes in U This is exactly the result we would get by equation (26.1) if we used the fact that U is known to alter the experiment, thinking of U as the new sample space, and counting only those outcomes in E that are also in U . There is nothing in equation (26.3) that is not already in equation (26.2), but it makes explicit the cancellations carried out in the last example to compute a conditional probability.

EXAMPLE 26.25

Four coins are tossed. Two of them fall within view of an observer, who sees that they are both heads. What is the probability, to this observer, that exactly three of the coins come up heads?

26.7 Conditional Probability

1125

The experiment is to toss four coins. An outcome consists of a string abcd of four letters, each of which is either an H or a T . The event E we are interested in consists of all such strings with exactly three H
s. E has four outcomes in it: E = THHH HTHH HHTH HHHT Let U be the event that at least two heads have come up, since this is what the observer knows when he or she sees two heads. Then U consists of all of the outcomes with exactly two heads, or exactly three heads, or exactly four heads. The number of these outcomes is 4 C2 + 4 C3 + 4 C4

= 6 + 4 + 1 = 11

Further, every outcome in E is also in U , so the number of outcomes common to E and U is four, the number of outcomes in E. By equation (26.3), the conditional probability that E happens, knowing that U happens, is Pr
E  U =

number of outcomes common to E and U 4 =  number of outcomes in U 11

The knowledge that at least two heads must come up makes the outcome of exactly three heads more likely than it would be without this information.

SECTION 26.7

PROBLEMS

1. Two coins are flipped. (a) What is the probability that the first one came up heads? (b) What is the probability that the first one came up heads, if we know that at least one came up heads? 2. Four coins are flipped. (a) What is the probability that exactly three came up heads? (b) What is the probability that exactly three come up heads, if we saw that one came up tails? 3. Four coins are flipped. (a) What is the probability that at least three came up tails? (b) What is the probability that at least three came up tails, if we saw two coins land tails? 4. Six cards are dealt from a standard deck without regard to order. (a) What is the probability that exactly two face cards were dealt? (b) What is the probability that exactly two face cards were dealt, if we know that a king was dealt?

5. What is the probability that four rolls of the dice total exactly 19? What is the probability of totaling exactly 19 if we know that one die came up 1? 6. What is the probability that two rolled dice sum to at least 9, if we know that one die came up even? 7. What is the probability that a five-card poker hand (unordered) has four aces, if we know that a four of spades was dealt? 8. What is the probability that a toss of seven coins will produce at least five heads, if we know that four of them came up heads? 9. Calculate the probability that four dice will all come up odd. What is this probability if we know that one came up 1 and another came up 5? What is the probability if we know that the second die came up 6? 10. Suppose A and B are events. Denote A ∪ B as the event containing all outcomes in A and all those in B. For example, if A = a b c d and B = c d e f, then A ∪ B = a b c d e f. Show that Pr
A ∪ B = Pr
A + Pr
B − Pr
A ∩ B Conclude that, if A and B have no outcomes in common, then Pr
A ∪ B = Pr
A + Pr
B. However, give an example to show that this equality fails to hold if A and B have outcomes in common.

CHAPTER 26

1126

26.8

Counting and Probability

Independent Events Informally, two events are independent if the knowledge that one occurs implies no information about the probability of the other occurring. This can be stated in terms of conditional probability.

DEFINITION 26.3

Events E and U are independent if Pr
E  U = Pr
E or Pr
U  E = Pr
U If E and U are not independent, then we say that they are dependent.

EXAMPLE 26.26

Consider the simple experiment of flipping one coin once. The sample space is S = H T, and the nonempty events are H, T and H T. We claim that the events E = H and U = T are dependent. The reason for this is that E ∩ U = ∅, so Pr
E ∩ U = 0, forcing Pr
E  U = 0 (equation (26.2)). However, Pr
E = 1/2, so Pr
E  U = Pr
E Similarly, U ∩ E = ∅, so Pr
U  E = 0, but Pr
U = 1/2, so Pr
U  E = Pr
U Therefore U and E are not independent. Although we have formally verified that E and U are not independent, the dependence of these events in this experiment is obvious. In one coin flip, if we know that one of the events H or T occurs, then we know that the other does not.

EXAMPLE 26.27

Draw two cards from a deck by first drawing one card, then drawing a second card (without replacement). The outcomes are pairs
a b, in which a is the first card drawn, and b the second. Order makes a difference here, since, for example, if a king of hearts is drawn first, then b could not be a king of hearts. The card drawn first influences the possibilities for the card drawn second. The sample space S consists of all pairs
a b with a and b cards, but a = b. There are 52 · 51 = 2652 outcomes in S. Consider the events E  a king is drawn first, and U  a jack, queen, king or ace is drawn second.

26.8 Independent Events

1127

Events in E are pairs
a b with a = king and b any card different from that drawn for a. There are four ways of drawing a king, and after this, 51 ways of drawing the second card, so there are 4 · 51 = 204 events in E. Then Pr
E =

204  2652

How many outcomes are in U ? U will consist of all pairs
a b, in which a and b are both drawn from the jacks, queens, kings and aces, together with all pairs
a b with only b drawn from the jacks, queens, kings or aces. In the first category, there are 16 possibilities for a and after that 15 for b, for 15 · 16 = 240 possible outcomes. In the second category with only b a jack, queen, king or ace, there are 36 choices for a and 16 for b, for a total of 36 · 16 = 576 outcomes. Therefore U has 240 + 576 = 816 outcomes. To compute the conditional probability Pr
E  U, we need to look at E ∩ U . This event consists of all outcomes
a b with a a king, and b drawn from the jacks, queens, kings and aces. There are four choices for a, and then 15 for b, so E ∩ U has 4 · 15 = 60 outcomes in it. Then number of outcomes common to E and U number of outcomes in U 60 =  816

Pr
E  U =

Since 204/2652 = 60/816, then Pr
E = Pr
E  U Similarly, it is routine to check that Pr
U = 816/2652 and Pr
U  E = 60/204, so Pr
U = Pr
U  E This means that E and U are not independent, and are therefore dependent. This makes sense intuitively. If we know U , then we know that a jack, queen, king, or ace was drawn second. This changes the probability that a king was drawn first.

EXAMPLE 26.28

We will repeat the experiment of the preceding example, but this time with replacement. That is, draw a card, record the result, then replace the card in the deck, shuffle, and draw again. Now the sample space S consists of all pairs
a b with a and b possibly the same card. There are 52 · 52 = 2 704 outcomes in S. Let E and U be the events of the preceding example. Now E consists of outcomes
a b, with a a king (four possibilities), but b can be any card (52 possibilities). There are now 4 · 52 = 208 outcomes in E, different from the preceding example because this experiment is different. Now Pr
E =

208  2704

This happens to be the same probability of this event in the previous experiment. Now there are more outcomes, but also more ways E can happen. Now consider pairs
a b in U . There are 52 possible choices for a, and 16 for b, for a total of 52 · 16 = 832 outcomes in U .

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Finally, the outcomes common to E and U are the outcomes
a b with a a king and b a jack, queen, king or ace. There are four ways to choose a, and (with replacement) 16 cards from which to draw b, for 4 · 16 = 64 outcomes common to E and U . Therefore number of outcomes common to E and U number of outcomes in U 64 = Pr
E = 832 The probability of E occurring is the same as the probability of E occurring, knowing that U occurs. Events E and U are therefore independent. Again, this makes sense intuitively. Because we are drawing with replacement, knowing that a face card was drawn second does not tell us anything about the first card drawn. Pr
E  U =

26.8.1 The Product Rule Equation (26.2) can be written Pr
E ∩ U = Pr
E  U Pr
U

(26.4)

This equation is called the product rule, and it is particularly useful in the context of independent events. If E and U are independent, then Pr
E  U = Pr
E or Pr
U  E = Pr
U. If Pr
E  U = Pr
E, then the product rule becomes Pr
E ∩ U = Pr
E Pr
U If Pr
U  E = Pr
U, we can use the fact that E ∩ U = U ∩ E and use the product rule to write Pr
E ∩ U = Pr
U ∩ E = Pr
U  E Pr
E = Pr
U Pr
E In either event, then, if E and U are independent, then Pr
E ∩ U = Pr
E Pr
U Conversely, suppose E and U are events, and Pr
E ∩U = Pr
E Pr
U. Then by the product rule, Pr
E ∩ U = Pr
E Pr
U = Pr
U Pr
E  U so Pr
E  U = Pr
E and therefore E and U are independent. We can summarize these conclusions as follows. THEOREM 26.3

Events E and U are independent if and only if Pr
E ∩ U = Pr
E Pr
U

(26.5)

This result is used in two ways. First, it provides another way of determining whether or not two two events are independent. Equation (26.5) is sometimes easier to check than the two conditions in the definition of independence. Second, equation (26.5) provides an often convenient way of computing the probability of an event E ∩U , consisting of outcomes common to E and U , as a product of the individual probabilities of E and U .

26.8 Independent Events

1129

Keep in mind, however, that the product rule applies only when E and U are independent. In this case, equation (26.5) reads “the probability of E and U is the product of the probability of E with the probability of U .

EXAMPLE 26.29

Suppose the probability of having a boy is 049, and that of having a girl, 051. A family has four children. What is the probability that exactly three of them are boys? Here the experiment is to have four children. An outcome is a string of four letters, for example, gbbg for girl, boy, boy, girl. The event E that we are interested in is that the family has three boys and one girl: E = bbbg bbgb bgbb gbbb The probability of any one of these is the product of the probabilities of each letter, g or b. The reason for this is that, in looking at four births !, any of these letters can be a g or a b independent of what the others are. Therefore Pr
bbbg =
049
049
049
051 = 006 = Pr
bbgb = Pr
bgbb = Pr
gbbb The probability of E is the sum of the probabilities of each outcome in E: Pr
E = 006 + 006 + 006 + 006 = 024 This probability is based on the information that girls are slightly more probable than boys in individual births. If boys and girls were equally likely (each probability 1/2), then we would proceed as follows. E has four outcomes in it, and the sample space has 24 = 16 outcomes. Now the probability of E is 4/16 = 025.

SECTION 26.8

PROBLEMS

1. Four coins are flipped. E is the event that exactly one coin comes up heads. U is the event that at least three coins come up tails. Determine whether E and U are independent. 2. Two cards are drawn from a standard deck, without replacement. E is the event that both cards are aces. U is the event that one card is a diamond and the other is a spade. Determine whether these events are independent.

comes up heads. U is the event that at least one die comes up 6. Determine whether these events are independent. 6. An experiment consists of picking two cards from a standard deck, with replacement. E is the event that the first card drawn was a king. U is the event that the second card drawn was an ace. Determine whether E and U are independent.

3. Two dice are rolled. E is the event that the dice total more than 11. U is the event that at least one die comes up an even number. Determine whether these events are independent.

7. Two cards are dealt from an honest deck, without replacement. E is the event that the first card dealt was a jack of diamonds. U is the event that the second card was a club or spade. Determine whether these events are independent.

4. A family has six children. E is the event that at least three are girls. U is the event that at least two are girls. Determine whether these events are independent.

8. Four coins are flipped. E is the event that the first coin comes up heads. U is the event that the last coin comes up tails. Determine whether these events are independent.

5. An experiment consists of flipping two coins and then rolling two dice. E is the event that at least one coin

9. Suppose a coin has been shaved so that the probability of tossing a head is 040. This coin is flipped

CHAPTER 26

1130

Counting and Probability

four times. What is the probability of getting at least two heads. What is the probability of getting exactly two heads? 10. A dishonest die comes up only 1, 4 or 6. This die is rolled three times. What is the probability that the total is 5?

26.9

11. A jar contains twenty marbles, eight red, eight blue and four green. The probability of drawing a blue marble is twice that of drawing a green, and three times that of drawing a red. What is the probability of drawing exactly two red marbles if three marbles are drawn at random, without replacement?

Tree Diagrams in Computing Probabilities Tree diagrams, in conjunction with the product rule, often provide a convenient way of computing certain probabilities in which the experiment can be broken down into a sequence of steps.

EXAMPLE 26.30

A cabinet has two drawers, labeled left and right. In the left drawer are three envelopes, e1 , e2 , and e3 , and in the right drawer are four vouchers, v1 through v4 . A drawer is chosen at random, then an object is selected at random from that drawer. What is the probability that the object drawn is v3 ? We may think of this experiment as proceeding in two stages, each of which is an experiment in its own right: Experiment 1—choose a drawer. Experiment 2—choose an object from a drawer. The tree diagram of Figure 26.5 displays the sequence of experiments. First, there is a probability of 1/2 of choosing either drawer, hence the numbers on the first two branches of the tree. There is a probability of 1/3 of choosing any of the envelopes, if drawer one was chosen. There is a probability of 1/4 of choosing any of the vouchers, if drawer two was chosen. The branches from each drawer to the last outcomes are conditional probabilities, since, for example, we can only pick an envelope if drawer one was chosen. This is where the product rule for probabilities comes in. To find the probability of choosing voucher v3 first follow the lower branch to drawer two, with probability 1/2, then the branch from drawer two to v3 , which has conditional probability 1/4. The product of these two numbers gives the probability of choosing v3 : Pr
v3  = Pr
drawer two Pr
v3  knowing drawer two was selected 1 1 1 · =  2 4 8 As we might expect by symmetry, the probability of choosing any envelope is the same (namely 21 · 13 = 16 ), and the probability of choosing any voucher is the same ( 21 · 41 = 18 ). Notice that the probabilities of all of the outcomes sum to one, as they must:     1 1 3 +4 = 1 6 8 =

The tree diagram is actually taking us through a sequence of conditional probabilities. The first stage of the experiment is to choose a drawer, but after this, the choice of envelope or voucher is a conditional probability, based on the choice of the drawer. In this example, there

26.9 Tree Diagrams in Computing Probabilities

1131

e1 (Prob = 1/6) 1/3 1/3 Left drawer

e2 (1/6)

1/3 e3 (1/6)

1/2

υ1

1/2 1/4 Right drawer

1/4

(1/8)

υ2 (1/8)

1/4

υ3 (1/8)

1/4

υ4 (1/8) FIGURE 26.5 Probabilities of outcomes of choosing a drawer, then an envelope or voucher.

were only two stages to the experiment, but more stages could be accommodated by continuing the tree with more edges from each stage to the next one. For the probability of any final outcome o, first identify each path leading to o. On each such path, multiply the probabilities on each branch of the path, and then add these results.

EXAMPLE 26.31

A room contains three urns. One urn has two baskets. One basket holds an envelope containing a $1 bill, an envelope containing $50 and an envelope containing $100. The other basket has three envelopes, one containing a quarter, one a half dollar and one a silver dollar. The second urn has two baskets. One basket holds an envelope containing nothing and an envelope containing $1 000, while the second basket has three envelopes, containing $50, $1, and $500. The third urn contains four baskets. The first basket has four envelopes, three holding nothing and the fourth, $200. The second basket has two envelopes, one with nothing and one with $300. The third basket has three envelopes, with, respectively, $1, $50, and $500. The fourth basket contains two envelopes, with $100 and $500. A person gets to choose one urn, then from it, one basket, and from it, one envelope. Determine the outcomes and the probability of each. This is an example of a problem that takes longer to state than to solve. Figure 26.6 shows a tree diagram for this experiment. By following paths to each payoff, as many times as it occurs as an end result, we obtain the following probabilities:

1132

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Counting and Probability

$1 $50

1/3 1/3

$100

1/3 $ 0.25

Basket 1 1/3 1/3

$ 0.50

Basket 2

1/2

1/3 $1 1/2

Urn 1

$0 1/2 1/3

$ 1000

1/2 Basket 1

$ 0.50 1/3

1/2

1/3 1/3

1/2

1/3 Basket 2

Urn 2

$ 500

$0 $0 $0 1/4 1/4

1/3

1/4 1/4

$ 200

Basket 1 1/2 1/2

1/4

$0 $ 300

1/4 Basket 2

$1

1/3

1/4

Urn 3

$ 50

1/3 1/4

1/3

Basket 3

$ 500 1/2 $ 100

Basket 4 1/2

$ 500 FIGURE 26.6

$1

Outcomes of choosing an urn, then a basket, then an envelope.

26.9 Tree Diagrams in Computing Probabilities

1133

  111 111 3 111 +3 + =  Pr
$0 = 322 344 3 4 2 16 Pr
$025 = Pr
$050 = Pr
$1 = Pr
$50 = Pr
$100 = Pr
$200 = Pr
$300 = Pr
$500 = Pr
$1 000 =

111 1 =  3 2 3 18 111 111 1 + =  323 323 9 111 111 111 111 7 + + + =  3 2 3 3 2 3 3 2 3 3 4 3 36 111 111 1 + =  3 2 3 3 4 3 12 7 111 111 + =  3 2 3 3 4 2 72 111 1 =  3 4 4 48 1 111 =  3 4 2 24 111 111 111 1 + + =  323 343 342 8 111 1 =  3 2 2 12

Again, these probabilities of all the outcomes must sum to 1.

SECTION 26.9

PROBLEMS

1. A cabinet has four drawers. Two of the drawers each contain two envelopes, each containing $10, one drawer contains one envelope with $5 and one envelope with $50, and one drawer contains one empty envelope and one envelope with $1 200. A person can choose any drawer, and from that drawer, any one envelope. What are the outcomes and their probabilities? 2. A room has three urns in it. One urn has two compartments, one with nothing in it, and one with a key which can open any of four safes. One safe has diamonds, one has stocks, one has cash, and one has a Cracker Jacks whistle of sentimental value to the owner of the urns. The other two urns each have three compartments. One of these urns has two empty compartments and one filled with Confederate currency. The third urn has one compartment filled with expensive perfume, one filled with stock certificates, and the third with the deed to a mansion. A person can pick an urn and

any compartment and, if the key comes up, any one of the safes. What are the outcomes and their probabilities? 3. A wealthy sultan has six automobile sheds which look identical from the outside. Each shed contains a number of identical containers, with each container closed, but holding one vehicle. One shed has two identical Fords and a Chevrolet (in their containers). A second shed has a VW Beetle (circa 1952) and a Porsche. A third shed has a Lamborghini and a very nice tricycle. A fourth shed has two Mercedes SUVs and a Honda Civic. A fifth shed has a World War I tank (partially destroyed in battle) and a Porsche. And the sixth shed has three mountain bicycles and a mint condition Stanley Steamer. A person can pick a shed, and then any container in that shed. What are the outcomes and their probabilities? 4. A traveler can choose to fly in a Piper Cub that seats one passenger,a company jet seating eight passengers,

1134

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Counting and Probability

or a jumbo jet seating seven hundred passengers. The traveler can pick any plane and any passenger seat on that plane. What is the probability that the traveler picks an odd numbered seat on the jumbo jet? 5. A person can choose any of five houses, each of which has four upstairs bedrooms. In one of the houses, three bedrooms are empty and one contains an antique chair worth $50 000. In another house, two bedrooms are

26.10

empty and one contains $1 000, and one contains a newly minted nickel. In a third house, each bedroom contains $500. And in the last two houses, two bedrooms contain $1 500 each, one contains $20, and one contains a person-eating lion which has not been fed recently. If the person is to pick a house and then a bedroom, what are the outcomes and their probabilities?

Bayes’ Theorem Bayes’ Theorem is named for the Reverend Thomas Bayes (1702–1761), although it only appeared in print after his death. It enables us to determine the conditional probability of E, knowing U , if we know the individual probabilities of E and U ∩ E (which is the event U and E ) as well as the conditional probabilities of U knowing E and U knowing E C . The rule is as close as we normally come in mathematics to “getting something for nothing,” and its derivation appears at first to be just a simplistic sequence of substitutions. If done in the right order, however, a very useful formula results. Begin with an experiment, and consider events E and U . Recall that E ∩ U = U ∩ E. Now Probability of E, Assuming U = Pr
E  U =

Pr
E ∩ U Pr
U

=

Pr
U ∩ E Pr
U

=

Pr
U  E Pr
E  Pr
U

(26.6)

in which we used the product rule (equation (26.4)) in the numerator to go from the next to last line to the last line. Now consider Figure 26.7, in which U and E are shown as typical sets of outcomes (events). We can split the outcomes in U into two kinds: those that are also in E, hence in U ∩ E, and those that are not in E, hence are in U ∩ E C . Furthermore, U ∩ E and U ∩ E C have no outcomes in common. Therefore Pr
U = Pr
U ∩ E + Pr
U ∩ E C  Now apply the product rule (26.4) to each of the probabilities on the right to obtain Pr
U = Pr
U  E Pr
E + Pr
U  E C  Pr
E C  Substitute this result into the denominator on the right in equation (26.6) to get Pr
E  U =

Pr
U  E Pr
E  Pr
U  E Pr
E + Pr
U  E C  Pr
E C 

(26.7)

Equation (26.7) is known as Bayes’ Theorem, although it is actually a special case of a more general result to be stated shortly (equation (26.8)). The theorem enables us to determine the conditional probability of E, knowing U , if we know: the probability of E, the probability of E C , the probability of U knowing E, and the probability of U knowing E C . The tree diagram in Figure 26.8 shows points representing the events E and E C , which together contain all the outcomes. Outcome U is given. From point E (the case that E occurs),

26.10 Bayes’ Theorem

1135

E

U U∩ E

U∩ EC

FIGURE 26.7 U consists of outcomes in U and in E
U ∩ E and outcomes in U that are not in E
U ∩ E C .

the two conditional events, U  E and U C  E are shown. And from vertex E C , the case that E does not occur, the two conditional events, U  E C and U C  E C are shown. Four of the branches are labeled with probabilities. The numerator in Bayes’ Theorem is the product of the probabilities on branches 1 and 2. The denominator is the sum of this product, and the product of the probabilities on branches 3 and 4. The paths ending in conditional probabilities involving U C are not relevant in computing Pr
E  U in this way. Bayes’ Theorem is all about drawing probability inferences from certain kinds of available information. Here is a typical application of the theorem. EXAMPLE 26.32

A factory produces wombles. Suppose, over a given day, the probability of producing a defective womble is 004. Suppose the probability that a defective womble will result in injury to a womble user is 002, while the probability that a womble user is injured through no fault in the product is 006. A lawsuit is in progress, in which the plaintiff has been injured and wishes a settlement from the company. The defense attorney wants some measure of whether a womble was to blame. The defense therefore wants to know the probability that the product was defective, if it is known that an injury occurred. Let U be the event that injury occurs, and let E be the event that a defective womble was produced. From the given information, Pr
E = 004

Pr
U  E = 002

and

Pr
U  E C  = 006

We want to determine Pr
E  U. Figure 26.9 shows the three probabilities just listed, together with the probability Pr
E C  = 1 − 004 = 096 which we infer from the principle of complementarity. We now have all the terms occurring on the right side of Bayes’ Theorem, so Pr
E  U = =

Pr
U  E Pr
E Pr
U  E Pr
E + Pr
U  E C  Pr
E C 
002
004 ≈ 00137
002
004 +
006
096

CHAPTER 26

Counting and Probability

Pr(U

/ E)

U/E

2 U/E

E

02

0.

UC/E

4

1

0.0

(E )

E Pr

Pr

(E

C

)

4 6

0.9

3

U/E C )

U/E C

Pr(

1136

EC

EC

0.0

6

UC/EC FIGURE 26.8

(simplest case).

Bayes’ theorem

U/EC FIGURE 26.9 Computing probabilities of a defective womble, knowing that an injury occurred.

In Figure 26.8, points for U C  E and U C  E C were included for completeness, although paths to these two points were not used in Bayes’ Theorem. In Figure 26.9, only the points representing events used in Bayes’ Theorem are included. There is a more general form of Bayes’ Theorem in which the single event E is replaced by k events E1      Ek . Now we find that Pr
Ej  U =

Pr
Ej  Pr
U  Ej  Pr
E1  Pr
U  E1  + Pr
E2  Pr
U  E2  + · · · + Pr
Ek  Pr
U  Ek 

(26.8)

for j = 1     k. Figure 26.10 illustrates this conclusion. Like Figure 26.9, this tree diagram contains only paths needed for terms in equation (26.8).

EXAMPLE 26.33

An airplane manufacturing company receives shipments of parts from five companies. The following table gives the percent of the total parts needs filled by each company, and the probability for each company that a part is defective (data taken over some period of time). A defective part is found. What is the probability that it came from Company 4?

26.10 Bayes’ Theorem

1137

U/E1

E1 U/E2 E2 U/Ek

Ek

C

E1

C

U/E1 C

E2

C

U/E2 C

Ek

U/EkC

FIGURE 26.10 General form of Bayes’

theorem.

Company 1 2 3 4 5

Supplied 15 2 4 48 27

Probability of Defect .03 .04 .01 .02 .08

Imagine a part has been chosen from the parts inventory. Let Ej be the event that the part came from company j, for j = 1 2 3 4 5. Let U be the event that the part is defective. We want to compute Pr
E4  U. From equation (26.8) we have enough information to do this: Pr
E4  U = =

Pr
E4  Pr
U  E4  Pr
E1  Pr
U  E1  + Pr
E2  Pr
U  E2  + · · · + Pr
E5  Pr
U  E5 
048
002
015
003 +
002
004 +
004
004 +
048
002 +
027
008

= 025

1138

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Counting and Probability

rounded to two decimal places. Notice that the percent supplied by each company (Column 2 of the table) is converted to decimals in thinking of these percents as probabilities. Just for comparison, consider the probability that the defective part came from the next largest supplier, Company 5. The denominator in applying equation (26.8) is the same as that used for Company 4, and we need only recompute the numerator: Pr
E5  U =


027
008
015
003 +
002
004 +
004
004 +
048
002 +
027
008

= 057 to two decimal places. The disparity is due to the fact that, while Company 5 supplies a little more than half the parts Company 4 does, it has four times the failure rate on its parts.

SECTION 26.10

PROBLEMS

1. A company hires people in four different groupings to install automobile carpet. The company has found that people with at least ten years experience generally produce products that are defective 1% of the time; people with five to ten years of experience, 3% of the time; people with one to five years experience, 57% of the time; and people with less than one year, 113% of the time. Of the total workforce preparing carpet, those with ten or more years comprise 45% of the total; those with five to ten years, 37%; those with one to five years, 7%; and those with less than one year, 11%. For each of these categories, calculate the probability that a defective installation was done by a worker in that group. 2. Suppose 37% of a company’s computer keyboards are manufactured in Los Angeles, and 63% in Detroit. Suppose 2% of the keyboards made in Los Angeles have some problem, and 53% of those made in Detroit are defective. Calculate the probability that a given defective keyboard was manufactured in Los Angeles, and also the probability that it came from Detroit. 3. A drug trial performed with seriously ill patients produces the following data: Adult Adult Boys Girls Men Women Survived at least two years 257 320 104 152 Survived less than one year 619 471 51 38

(a) If a patient is selected randomly from this pool, what is the probability of choosing an adult male, if we are told that the person survived at least two more years?

(b) What is the probability of choosing a girl if we are told that the person survived less than one year? 4. A certain herb is grown in three places on Earth. It is known that this herb has medicinal value in treating certain illnesses, but that sometimes it is fatal. In the accompanying table, column two gives the percent of the total herb production of the world that is grown in that location, and the third column gives a percent of the herb from that location that is fatal.

Lower Amazon Sub-Saharan Africa Newark, New Jersey

% of Total 62 25 13

Percent Fatal 14 6 1/2

For each location, calculate the probability that a sample of the herb came from that location, if it was found to be fatal. 5. Target pistols are manufactured in six cities. The following table gives the percent of the total production that comes from each city, together with the probability that the gun explodes upon firing. City City City City City City

1 2 3 4 5 6

% of Total 15 7 21 4 9 44

Probability of Exploding .02 .01 .06 .02 .03 .09

Suppose a gun explodes. Calculate, for each city, the probability that the gun came from there.

26.11 Expected Value

26.11

1139

Expected Value For anyone visiting a carnival or gambling casino, the expected value of a game or experiment is a very useful number to be able to compute. Here is an example to illustrate the idea.

EXAMPLE 26.34

A carnival game has a large wheel with the numbers 1 through 50 equally spaced around the outer rim. It costs $1 to bet on which number will end up at the top when the wheel is spun and allowed to come to rest. A winner receives $25. Is this a profitable game to play? In view of the payoff, which is twenty-five times the cost of making a bet, it apparently is believed that such a game is a winning proposition for the player. At least, the game is popular and seems to draw a lot of players. On the other hand, we should be suspicious, since the carnival is a profit-making enterprise and the owners of the game must be making money or they would change the game. Who is right? We will answer this question by trying to define a sense of how much we should expect to win or lose each time we play. On one spin of the wheel, there are fifty equally likely outcomes, assuming an honest wheel (a reasonable assumption, since the player can bet on any number). On any number we pick, the probability of winning is 1/50, and the probability of losing is 49/50. Since winning pays $25 and losing costs $1, compute 49 1
$25 −
$1 = −$048 50 50 This is the expected value of this game from the player’s perspective. A player should expect on average to lose (because of the negative sign) forty-eight cents on each spin. Therefore the "house" expects to win this much on each spin, with each player. Although the game owner will occasionally make payouts, the owner stands on average to make nearly a half dollar per player per game. Thought of in a slightly different way, the player will on average win $25 , for 1/50, or 2% of the time, and lose $1 for 49/50, or 98% of the time. Although the winning payoff is many times the loss on each game, a person will most likely lose so many more times than he or she wins that on average each person will contribute forty-eight cents to the carnival for each game. In general, we obtain the expected value of an experiment by multiplying the probability of each outcome by the value of that outcome, with a plus sign for value gained, and a negative sign for value lost, and adding the results. Thus, the expected value has a perspective, since a loss for one person may be a win for another. Of course, the concept of expected value assumes that each outcome has been given a value or payoff of some kind.

EXAMPLE 26.35

Twelve coins are flipped. We win $8 every time at least eight heads come up, and lose $2 otherwise, with no significance placed on the order of the coins. What is the expected value of this game?

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Counting and Probability

The number of ways of obtaining at least eight heads is the sum of the number of ways of obtaining exactly eight, nine, ten, eleven or twelve heads. This number is           12 12 12 12 12 + + + + 8 9 10 11 12 = 495 + 220 + 66 + 12 + 1 = 794 Since there are 2

12

outcomes, the probability of obtaining eight or more heads is Pr
eight or more heads =

794  4096

The probability of this not occurring is Pr
fewer than eight heads = 1 −

794 3302 =  4096 4096

The expected value of this game is 794 3302 252
8 −
2 = −  4096 4096 4096 or about −00615 dollars. On average, the owner expects to make slightly over six cents on each game.

EXAMPLE 26.36

Based on performance over the past thirty years, a salesman estimates the following probabilities for commissions on sales over the coming year: 025 probability of getting $400 from shoe sales, 015 probability of getting $200 from trouser sales, 035 probability of getting $600 from sport jacket sales, 020 probability of getting $1 000 from sales of complete golf outfits, and 005 probability of getting $300 from sales of derby hats. He wants to calculate his expected commission to estimate his income for tax purposes. The information is most easily read by making a table: Commission Probability

400 0.25

200 0.15

600 0.35

1000 0.20

300 0.05

Notice that the probabilities add up to 1, since presumably all possibilities for commission income have been included. The expected commission income is 400
025 + 200
015 + 600
035 + 1000
020 + 300
005 = 555 dollars.

SECTION 26.11

PROBLEMS

1. A game consists of flipping seven coins, without regard to order. If at least three heads come up, the player wins $5. If not, the player loses $9. What is the player’s expected value for this game?

3. Six cards are dealt from a standard deck, without regard to order. If one or more aces comes up, the player wins $45. Otherwise he or she loses $30. What is the player’s expected value?

2. Five dice are tossed. If at least three come up 5 or 6, the house pays $50. If not, the house wins $4. What is the player’s expected value?

4. Nine coins are flipped. If seven or more come up tails, the player wins $50. Otherwise the player loses $3. What is the player’s expected value?

26.11 Expected Value

1141

5. Twenty marbles, labeled 1 through 20, are in a jar. A game consists of reaching in blind and pulling out three marbles. If at least two are even, the player wins $3. Otherwise, the player loses $7. What is the player’s expected value?

7. Seven hats are in a closet. Four are black and three are red. A person reaches in and pulls out four hats, without regard to order. If at least two are red, the player wins $10. Otherwise the player loses $5. What is the player’s expected value?

6. Sixteen dice are rolled, and unordered results are recorded. If at least five sixes come up, the player wins $70, otherwise the player pays $4. What is the player’s expected value?

8. Six coins are flipped, then four dice are rolled. If four heads or a total of at least 24 on the dice result, the player wins $25. Otherwise the player loses $3. What is the expected value of this game?

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CHAPTER

27

BELL CURVES, NORMALLY DISTRIBUTED DATA, BINOMIAL DISTRIBUTION, STUDENT t-DISTRIBUTION, CORRELATION COEFFICIENT, CONFIDENCE INTERVAL, REGRESSION LINE

Statistics

With this background on counting and probability, we can begin to look at some of the ideas behind statistics.

27.1

Measures of Center and Variation Many processes generate numbers, in tables of data, graphs, or some other format. Such numbers are called data. A statistic is a measure that assists in the analysis of data in order to extract information and draw conclusions from the data. In this section we will discuss two kinds of elementary statistics. The first kind are measures of center, attempting to quantify in some sense the “middle” or “center” of a data set. The second kind of measure attempts to quantify the “spread” of the data, giving a sense of how the numbers vary from each other.

27.1.1

Measures of Center

Suppose we have N numbers, A1 , A2 ,     AN , not necessarily all distinct.

DEFINITION 27.1

Mean or Average

The mean, or average, of numbers A1 , A2 ,     AN is their sum, divided by N : Mean of A1  A2      AN =

A 1 + A 2 + · · · + AN  N

This number is often denoted x. Sometimes the Greek letter  (mu) is also used. 1143

1144

CHAPTER 27

Statistics

EXAMPLE 27.1

The mean of the numbers 16, −25, 27, 12, 15 and 31 is x=

16 − 25 + 27 + 12 + 15 + 31 38 =  6 3

As this example shows, the mean of a set of numbers need not be one of the numbers in the set. In computing means, some numbers may occur more than once in the data list. For example, in a list of daily high temperatures in a city over a thirty year period, some temperatures may be repeated. These are still listed because they accurately record the data of interest. The mean of a data set is a statistic that provides what is in a sense a “balance point” for the data. Consider a standard 12 inch ruler, and suppose the data set consists of the numbers 2, 3, 35, 4, 8, 11 and 115. Imagine a one gram weight placed at each of these points on the ruler, as in Figure 27.1. At what point should a fulcrum be placed so that the ruler balances? The answer is, it should be placed at the mean of the data set, which is 2 + 3 + 35 + 4 + 8 + 11 + 115 = 61429, 7 just past the half-way mark on the ruler. In this way the mean behaves like the center of mass of the data. Mean is a useful statistic, but it must be treated with care. x=

6.1429

FIGURE 27.1

The mean is the balance point for the data set.

EXAMPLE 27.2

Last year there were eight graduates of a very prestigious electrical engineering progam. Upon graduation, they all received good job offers with starting salaries listed below: Person

Starting Salary

Jones Smith Connolly Andrews Finch Hatton Douglas Seagram

$80,000 $85,000 $78,000 $90,000 $80,000 $87,000 $92,000 $6,000,000

The average starting salary is x = $824000. This is very misleading. The first seven graduates listed certainly did very well, but nothing like this. What the data do not reveal is that Seagram hit 200 home runs in 250 at bats in college and could pitch every other day, either right handed

27.1 Measures of Center and Variation

1145

or left handed, with a 153 mile per hour fast ball and a curve ball that breaks at an angle of 3/4 radians. In this data set, Seagram would be called an outlier, which is a point very far removed from all of the others. Outliers skew the mean and invite false conclusions if care is not taken. In some analyses of (fairly large) data sets, the highest and lowest numbers are dropped before computing a mean. For example, in some figure skating scoring systems, the high and low scores are deleted and the other scores averaged. A different measure of center is given by the median. This number locates the “middle” of a data set, without taking into account the actual magnitude of the numbers themselves.

DEFINITION 27.2

Median

Suppose we are given an ordered data set A1 , A2 ,    , AN , in which Aj ≤ Aj+1 for j = 1     N − 1. The median of this set is a number M having the property that half of numbers in the data set are ≤ M, and half of the numbers are ≥ M.

The mean gives a center of gravity (weighted center) of the data set, taking the magnitude and sign of each number into account. By contrast, the median gives a geometric center (actual middle of the list), and only assumes that the numbers are written in nondecreasing order, without other regard to their magnitude. Like the mean, the median may or may not itself be a number in the data set. It is easy to verify that if: N is odd, then M = A
N +1/2 and N is even, then M =

AN/2 + A
N/2+1 2



In the case that N is odd, one number, namely A
N +1/2 , is at the exact center of the ordered list, so in this case the median is a number in the list. When N is even, the median is the average of entry N/2 and entry
N/2 + 1 in the list, counting from the left, and may or may not be a number in the list.

EXAMPLE 27.3

Find the median of the data set −4 16 1 2 2 3 7 5 15 −4 2 1 1 3 9 8. First order (or sort) the list: −4 −4 1 1 1 2 2 2 3 3 5 7 8 9 15 16 Here N = 16, so A 8 + A9 2+3 = = 25 2 2 In this case the median is not in the set. There are eight numbers in the data set to the left of 25, and eight numbers to the right: M=

−4 −4 1 1 1 2 2 2 M = 25 3 3 5 7 8 9 15 16 : ;< = : ;< =

1146

CHAPTER 27

Statistics

EXAMPLE 27.4

Consider the data set −4 16 1 2 2 3 7 5 15 −4 2 1 1 3 9 8 1 As an ordered list, this is −4 −4 1 1 1 1 2 2 2 3 3 5 7 8 9 15 16 with N = 17 numbers, an odd number. If we denote these numbers B1      B17 , the median is M = B
N +1/2 = B9 = 2 and we can write −4 −4 1 1 1 1 2 2 M = 2 3 3 5 7 8 9 15 16 ;< = : ;< = : The median is 2, but it is the third 2 from the left in the ordered list, not one of the other 2’s. In this case the median is a number in the list. Data is frequently given by a frequency table, particularly if N is large. This is a table that gives each number according to its frequency in the list, saving the trouble of writing the number down as many times as it is repeated.

EXAMPLE 27.5

Data (already ordered) is given in the frequency table Data Point Frequency

−3 5

−1 2

4 12

7 6

12 12

17 21

The data has −3 listed five times, −1 two times, and so on. The sum of the frequencies is N = 58, and this is the number of data points. The mean is
−3
5 +
−1
2 + 4
12 + 7
6 + 12
12 + 17
21 287 =  58 29 approximately 9897. Because the number of data points is even, the median is
A29 + A30 /2. Counting the frequencies from the left, the 29th number in the ordered list occurs in the group of twelve 12’s, and equals 12. The 30th number is also 12, so the median is 12. =

27.1.2 Measures of Variation Measures of variation are indicators of how “spread out” the data is, as opposed to where a “center” might be located. The simplest measure of variation is the range.

DEFINITION 27.3

Range

The range of a data set is the difference between the largest number in the set and the smallest number in the set.

27.1 Measures of Center and Variation

1147

EXAMPLE 27.6

In the data set −4 −4 1 1 1 1 2 2 2 3 3 5 7 8 9 15 16, the range is 16 −
−4 = 20. If these numbers are marked on the number line, the number at the left is twenty units from the number at the right of the list.

Note that the data set need not be ordered to compute the range. However, if the set is large, it is certainly easier to locate the largest and smallest numbers if the list is ordered. Another measure of variation is the standard deviation, which considers variations of data points from the mean of the data.

DEFINITION 27.4

Standard Deviation

Let A1      AN be numbers in a data set, having mean x. 1. If these numbers represent a sample within a (generally much) larger population, then the standard deviation of this sample is > ? N ? 1 
A − x2  s=@ N − 1 j=1 j 2. If these numbers represent the entire population, or total set of items under consideration, then the standard deviation is > ? N ?1   =@
A − x2  N j=1 j

For example, an x-ray unit might take data at 4000 points. The readings at these points would be the entire population of data, and  would be used for a standard deviation. It might be useful, however, to periodically take samples of 30 of these points, for which we would use the standard deviation s. The numbers s2 and  2 (squares of the standard deviation) are called the variance of the data, and are used in some contexts. The standard deviation of A1      AN is zero if and only if each Aj = x. In this case the data points are all the same and there is no deviation from the mean. However, if one or more Aj has a large difference from the mean, then the standard deviation will be correspondingly large. A large standard deviation indicates data that are “spread out”, while a small standard deviation corresponds to data that are clustered together, hence all lie fairly close to the mean. Of course, the idea of variation from the mean makes no sense if there is only one number, so N > 1 is assumed in defining the standard deviation s. It is a routine exercise in algebra to show that

s=

> 2  ? N @ N j=1 A2j − Nj=1 Aj N
N − 1



1148

CHAPTER 27

Statistics

Some statisticians use 1/N in place of 1/
N − 1 in defining s. We will use 1/
N − 1 because, for a given s, if x is known, then any N − 1 of the data points determine the remaining one.

EXAMPLE 27.7

Consider the data set consisting of the numbers 11, 13, 16, 2, 21, 22, 24 and 25. These numbers are all fairly close together, with a range of only 25 − 11 = 14. The mean of these numbers is x=

11 + 13 + 16 + 2 + 21 + 22 + 24 + 35 = 2025 8

For the standard deviation, first compute the sum of the squares of the differences of the numbers from the mean:
11 − 20252 +
13 − 20252 +
16 − 20252 +
2 − 20252
21 − 20252 +
22 − 20252 +
24 − 20252 +
25 − 20252 = 19650 The standard deviation is  s=

19650 = 05298 7

to four decimal places. This small standard deviation is due to the fact that these numbers are fairly close to one another with not much spread. Contrast this standard deviation with a similar calculation for the following data set, which is more “spread out.”

EXAMPLE 27.8

Consider the data set 11, 46, 92, 157 and 28. The range is 269. Routine calculations give x = 1172 and  = 94943. The large standard deviation reflects the fact that the numbers in the data set are spread out. Greater spread would result in a larger standard deviation. For example, if 11 were replaced by −127, making√the range 4070, then the mean of the new data set is 896, and the standard deviation is s = 48119/4 ≈ 10968. Figures 27.2 and 27.3 show scatter plots of the data in Examples 27.7 and 27.8, respectively, giving a visual sense of the connection between spread of the data and standard deviation. As another example of the idea of s as a spread, we would expect the standard deviation of the data set of electrical engineering starting salaries to be very large. A routine calculation shows that this standard deviation is approximately 2091400, all due to the outlier Seagram.

27.1 Measures of Center and Variation

1149

3.0

2.5

2.0

1.5

1.0

0.5

0 0

1

2

3

4

5

6

7

8

9

FIGURE 27.2 Data of Example 27.7.

30

25

20

15

10

5

0 0

1

2

3

4

5

6

FIGURE 27.3 Data of Example 27.8.

SECTION 27.1

PROBLEMS

1. Sometimes the range of a data set is used to make a rough estimate of the standard deviation, according to the rule s≈

1
largest number in the list − smallest number 4

Compute s, then try this estimate for the following data sets: (a) −5 −2 −15 0 0 1 1 1 2 7 9 116 (b) 7 1 −4 2 2 1 4 1 1 2 6 3 (c) 5 −2 −10 5 2 2 5 3 5

CHAPTER 27

1150

Statistics

2. Compute the mean, median, and standard deviation s for each of the data sets: (a) −4 −6 25 3 8 5 −3 −4 8 3 22

−3 4

−1 2

0 6

1 4

3 12

4 3

Find the mean, median, and standard deviation s for this data set.

(b) 1 1 1 −1 2 3 −1 4 2

4. Data are given by the following frequency table:

(c) 3 −4 2 15 −4 −4 2 1 7

Data point Frequency

(d) 93 95 97 10 84 87 88 88 41 (e) −16 −14 −10 0 0 1 1 3 5 7 3. Data are given by the following frequency table:

27.2

Data point Frequency

−12 4

−97 2

−8 6

−76 4

−51 12

4 3

Find the mean, median, and standard deviation s for this data set.

Random Variables and Probability Distributions A random variable on an experiment is a function that assigns a real number to each outcome of the experiment, based in some specified way on the random outcomes of the experiment. If X is such a random variable, then X
o is a real number for each outcome o of the experiment.

EXAMPLE 27.9

Consider the experiment of flipping three coins. Each outcome is a string of three letters, each an H or T . If o is an outcome, let X
o = number of heads in o. X assigns a number to each outcome. In this example, for any outcome o, X
o must be one of the numbers 0, 1, 2, or 3. The outcomes of this experiment are o1 = HHH

o2 = HHT

o3 = HTH

o4 = THH

o5 = HTT

o6 = THT

o7 = TTH

o8 = TTT

and X
o1  = 3

X
o2  = X
o3  = X
o4  = 2

X
o5  = X
o6  = X
o7  = 1

and

X
o8  = 0

For any repetition of this experiment, any outcome may occur, hence X is a random variable.

EXAMPLE 27.10

Consider the experiment of flipping six coins, but repeating these six flips fifteen times. An outcome consists of a list of fifteen strings of six letters each, each letter being an H or T . If o is an outcome, let X
o equal the total number of tails that came up in the fifteen repetitions of six flips. X
o can be any number from 0 through 90, inclusively and X is a random variable, because for any fifteen repetitions of the six flips, the outcome is random.

27.2 Random Variables and Probability Distributions

DEFINITION 27.5

1151

Probability Distribution

A probability distribution on a random variable X is a function P which assigns a probability P
x to each value x that the random variable can assume.

 If x denotes a summation over all values x that the random variable can assume, then a probability distribution must satisfy  P
x = 1 x

and 0 ≤ P
x ≤ 1 for each outcome x. Both of these requirements are consistent with our understanding of a probability function. The only new wrinkle is that now we have a probability function not on the outcomes themselves, but on the values of the random variable X. The probability P
X
o is the probability, not of the outcome o itself, but of X
o, the value of the random variable at o. We are using P for this probability distribution, to distinguish it from a probability Pr that may be defined directly on the outcomes of the experiment. Pr
o is the probability that the outcome o occurs, and P
x that the random variable assumes the value x. We could therefore write  is the probability  P
x as P
X
o, the first sum being over the values x that X takes on, and the second x o sum over the outcomes o of the experiment.

EXAMPLE 27.11

Flip three coins. There are 23 = 8 outcomes
o1      o8 , each with probability Pr
oj  = 1/8 Define a random variable X as follows. For each outcome oj , let X
oj  = number of tails in oj . Then X has the following values: X
HHH = 0

X
HHT = X
HTH = X
THH = 1

X
TTH = X
THT = X
HTT = 2

X
TTT = 3

There are four values this probability distribution can assume, namely 0, 1, 2, and 3. A probability distribution on X is given by assigning a probability to each value X can assume. In this example, assign 1 3 P
0 =  P
1 =  8 8 3 1 P
2 = and P
3 =  8 8 The probability function Pr acts on the outcomes of the experiment, and every outcome is equally likely. The probability distribution P acts on the values of the random variable, giving a probability for each value of the random variable to occur. Not all values that the random variable assumes are equally likely. Three of the eight outcomes outcomes have two tails, so P
2 = 3/8 and so on.

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Statistics

The mean and standard deviation for a random variable are defined as follows.

DEFINITION 27.6

Mean and Standard Deviation

Let X be a Random Variable. Suppose X has probability distribution P. 1. The mean of X is the number  defined by   = xP
x x

2. The standard deviation of X is denoted , and is defined by  =
x − 2 P
x The sum

x

 x

is over all values that the random variable can assume.

The expression for  is reminiscent of the expected value of an experiment or game. Each term of the sum for  is the product of a payoff x and the probability that the random variable achieves this value. As a special case, if the numerical values of X are x1 , x2 ,   , xN and if X is equally likely to achieve each of these values, then P
xj  = 1/N and the formula for  becomes =

N  j=1

 xj

1 N

 =

N 1  x = x = mean of the data set x1      xN  N j=1 j

Further, in this case 0 =

n 

0
xj − 2 /N =

j=1

N

j=1
xj

− x2

N



the familiar formula for the standard deviation of a population. It is a routine calculation to show that  =

x2 P
x − 2  x

EXAMPLE 27.12

If we flip four coins, there are 24 = 16 outcomes: o1 = HHHH

o2 = TTTT

o3 = THHH

o4 = HTHH

o5 = HHTH

o6 = HHHT

o7 = HTHT

o8 = HHTT

o9 = HTTH

o10 = THHT

o11 = THTH

o12 = TTHH

o13 = HTTT

o14 = THTT

o15 = TTHT

o16 = TTTH

27.2 Random Variables and Probability Distributions

1153

Define a random variable X by letting X
o equal the number of tails in the outcome o. From the list of outcomes, we have: X
o1  = 0

X
o2  = 4

X
o3  = X
o4  = X
o5  = X
o6  = 1 X
o7  = X
o8  = X
o9  = X
o10  = X
o11  = X
o12  = 2 and X
o13  = X
o14  = X
o15  = X
o16  = 3 X can take on five numerical values. We can assign a probability to each numerical value of X by defining P
0 = P
4 =

1  16

4 6 4  P
2 =  and P
3 =  16 16 16 For example, the value 3 (for three tails) occurs four times out of the sixteen outcomes of the experiment, so its probability of occurring as a numerical value of X is 4/16. Notice that these probabilities sum to 1. The mean of X is   = xP
x P
1 =

x



         1 1 4 6 4 =0 +4 +
1 +
2 +
3 = 2 16 16 16 16 16 On average, with four flips, we expect to see two tails. The standard deviation for X is   =
x − 2 P
x x



     1 1 4 +
4 − 22 +
1 − 22 16 16 16     6 4 +
2 − 22 +
3 − 22 = 1 16 16

=
0 − 22

Random variables are not restricted to have only finitely many outcomes. For example, if the experiment is to pick a positive integer k, and X
k = k, then X is a random variable which assumes infinitely many values. This leads to the following definition, which includes a subtlety involving orders of infinity.

DEFINITION 27.7

Random Variables

A discrete random variable is one that assumes either a finite or countable number of values. A continuous random variable is one that assumes an uncountable infinity of values.

CHAPTER 27

1154

Statistics

An infinite set is called countable if it can be put into one-to-one correspondence with the set of positive integers. For example, suppose E consists of all powers 2n , with n a positive integer. The correspondence n ↔ 2n , for n any positive integer, matches each positive integer with a unique number in E, and conversely. E is therefore countable, or countably infinite. An infinite set that cannot be put into one-to-one correspondence with the positive integers is said to be uncountable, or uncountably infinite. The real numbers form an uncountable set, because it can be shown that there is no one-to-one match between the positive integers and the real numbers. The real numbers have a higher order of infinity than the positive integers, even though both sets are infinite. As an example of a continuous random variable, suppose the experiment is to choose a number (not necessarily an integer) between 0 and 1 inclusive. There are uncountably many such numbers. If X
k = k for 0 ≤ k ≤ 1, then X is continuous because the number of values X can assume is uncountably infinite.

SECTION 27.2

PROBLEMS

1. Suppose two dice are rolled. If o is an outcome, let X
o = the sum of the numbers on the dice. Determine the probability distribution of X, as well as the mean and standard deviation for X. 2. Four coins are tossed. If o is an outcome, let  1 if two or more tails come up X
o = 3 otherwise. Determine the probability distribution of X, and its mean and standard deviation. 3. A wheel having the numbers 1 through 20 equally spaced about its circumference, is spun, and the number that is on top when the wheel comes to rest is recorded. If o is an outcome, let X
o = the number of factors in the prime factorization of o. For example, if 7 comes up, then X
7 = 1 because 7 is prime. But X
15 = 2 because 15 = 3
5, and X
12 = 3 because 12 = 2
2
3. By convention, 1 is not prime and has no prime factors. Determine the probability distribution of X, and the mean and standard deviation of X.

27.3

4. Suppose two dice are rolled. If o is an outcome, and the two dice show different numbers, let X
o be the quotient of the larger number divided by the smaller. If o is an outcome in which both dice come up the same, let X
o be this common value. Determine the probability distribution of X, and its mean and standard deviation. 5. Two cards are picked in succession (without replacement) from a standard deck. If o is an outcome, let X
o equal the sum of the numbers on the cards if both cards are numbered; X
o = 11 if exactly one card is a jack, queen, king or ace; and X
o = 12 if two of the cards are in the jack, queen, king, ace group. Determine the probability distribution for X, as well as its mean and standard deviation. 6. An experiment is to choose any integer from 1 to 30, inclusively. If o is an outcome, let X
o = 1 if o is divisible by 2, X
o = 2 if o is divisible by 3 but not 2, and X
o =  otherwise. Determine the probability distribution of X, and its mean and standard deviation.

The Binomial and Poisson Distributions Not every distribution that can be defined is interesting. However, the binomial and Poisson distributions have important applications. This section is devoted to these two distributions.

27.3.1 The Binomial Distribution The binomial distribution is used in a very specific setting. Suppose we have some experiment or procedure that involves a fixed number of trials (repetitions). For example, in flipping a coin

27.3 The Binomial and Poisson Distributions

1155

fifty times, the experiment consists of the fifty flips, and each flip is a trial. Assume that the following conditions are satisfied. 1. Each trial must have exactly two outcomes. These might be called A and B, or S (success) and F (failure), or H and T (for a coin), or some other designation. 2. The trials are independent. The outcome of one trial is unaffected by the outcome of any other trial. 3. For each trial, the probabilities remain constant. For example, with coin tosses, each toss might have 1/2 probability for a head, and 1/2 for a tail, and these probabilities do not change as more trials are performed. Now suppose N trials are carried out, and the result of each trial is S or F . Define X
k = the number of times S comes up exactly k times in the N trials. Then X is a random variable, taking on the values 0, 1 ,   , N . Suppose S has probability p on any one trial, meaning that F has probability q = 1 − p. Then the probability distribution P of X is the binomial distribution   N P
x = px q N −x for x = 0 1     N (27.1) x   N Recall that or N Cx is the number of ways of choosing x objects from N objects without x regard to order and is given by   N! N =  x x!
N − x! To see why equation (27.1) holds, consider the probability of getting S exactly x times in the  N trials. Imagine picking x boxes out of N boxes, without regard to order. There are N ways of doing this. For each such way, the probability of that way occurring is px q N −x , x because S comes up x times, each with probability p, and therefore F comes up N − x times, each with probability q. Therefore the probability of getting S exactly x times is P
x = (number of ways this outcome occurs) · (probability of each occurrence)   N = px q N −x  x Notice that this is the term having p to the x power in the binomial expansion
p + qN . Since p + q = 1,  N    N P
x = px q N −x =
p + qN = 1N = 1 x x x=0 as is required for the probability distribution of a random variable.

EXAMPLE 27.13

A district attorney is deciding on strategy for considering a plea bargain. The D.A. knows that jurors make a correct vote for guilt or innocence about 79% of the time. There are twelve jurors, who vote independently, and it takes a vote of at least ten to convict. What is the probability that at least ten jurors will reach a correct decision?

1156

CHAPTER 27

Statistics

We want the probability of at least ten out of twelve votes being correct (as decided by the evidence and presentation). Let N = 12, p = 079 and q = 1 − 079 = 021. Since at least ten jurors means exactly ten, exactly eleven or exactly twelve, compute       12 12 12 P
10 + P
11 + P
12 =
07910
0212 +
07911
021 +
07912 10 11 12 12 · 11
07910
0212 +
12
07911
021 +
07912 2

=

≈ 0523 This is the probability that at least ten jurors reach a correct decision. As another example, the probability that exactly eight jurors reach a correct decision is   12 P
8 =
0798
0214 ≈ 0146 8 For a binomial probability distribution, the mean and standard deviation have a particularly simple form.

THEOREM 27.1

For a binomial probability distribution on N trials,  = Np and  =



Npq

These are routine computations. For the mean, write     N N   N N k N −k = k pq = k pk q N −k k k k=0 k=1 =

N 



N

k=1

because

 k

N k



N −1 k−1



pk q N −k 


N − 1! kN ! =N k!
N − k!
k − 1!
N − k!   N −1 =N  k−1 =

Now change the summation index to m = k − 1 to write =

N  k=1

= Np

 N

N −1 k−1

N −1   k=0



N −1 m

k N −k

pq 

=

N −1  m=0

 N

N −1 m

 pm+1 q N −1−m

pm q N −1−m = Np
p + qN −1 = Np

because p + q = 1. A similar calculation verifies that  =

√

Npq.

27.3 The Binomial and Poisson Distributions

1157

EXAMPLE 27.14

An assembly line produces automatic veebles. It has been found that the probability that a veeble is defective (fails to veebulate) is 002. If 40 veebles are selected at random, what is the probability that no more than three are defective? Let N = 40 and p = 002, so q = 098. The random variable X is defined by X
x = number of veebles that are defective, for x = 0 1     40. The probability distribution is binomial, and P
x = probability that x of the forty veebles are defective. If no more than three are defective, then none, one, two or three could be defective. We want to know         40 40 40 40 0 40 39 2 38 P
0 + P
1 + P
2 + P
3 = pq + pq + pq + p3 q 37 0 1 2 3 =
09840 + 40
002
09839 + 20
39
0022
09838 + 20
13
38
0023
09837 ≈ 0992 This probability makes sense intuitively. The probability of one veeble being defective is quite small, so it should be very likely that fewer than four are defective. The mean of this probability distribution is  = Np = 40
002 = 08 The standard deviation is =

27.3.2



Npq =



40
002
098 ≈ 0885

The Poisson Distribution

Suppose we have some specific interval or segment in mind. This could be an interval in time, or, in terms of a space unit, some segment such as a length, distance, area, or something else. We are often interested in the occurrence of some event falling within this interval or segment. Let X
x = the number of times x occurs in the given interval. Then X is a random variable, assuming that the occurrence of falling or not falling in the interval is a random event. The probability that the event falls in the interval x times can be shown to have the form x e−  (27.2) x! in which  is a ratio of certain occurrences per interval, according to the setting of the problem. P is a probability distribution on the random variable X, and is called the Poisson distribution. The binomial distribution is discrete, and in fact finite, since it deals with the probability of some event occurring a certain number of times out of a given number. By contrast, the Poisson distribution is not finite, since it can take on values 0 1 2    . In the Poisson distribution, it is assumed that the occurrences under consideration are independent of each other, and are uniformly distributed over the interval of interest. The occurrences are therefore not allowed to cluster in one particular part of the interval. P
x =

EXAMPLE 27.15

A game consists of tossing darts onto a large flat mat that has been divided into 450 blocks of 6 square inches each. In one session, 370 darts were thrown. Suppose we want the probability that one block was hit exactly twice or exactly four times.

1158

CHAPTER 27

Statistics

We will use the Poisson distribution because the game involves an event (dart hitting the floor) occurring over an interval, which in this case is a 6-square-inch block of the floor. First, compute 370 = 0822 450 darts per block. Let P
x = x e− /x! for x = 0 1 2   . We want the probability of a block being hit exactly two or exactly four times. Thus, compute =


08222 e−0822
08224 e−0822 + ≈ 0157 2! 4! Since 450
0157 = 7065, we would expect on average that 71 blocks will be hit exactly two or four times. For contrast, suppose only 80 darts were thrown. Now take P
2 + P
4 =

=

80 = 0178 450

so
01782 e−0178
01784 e−0178 + = 00133 2! 4! Since 450
00133 = 5985, we would expect six of the blocks to be hit exactly two or four times if 80 darts are thrown. P
2 + P
4 =

It is possible to show that the mean of the Poisson distribution (27.2) is , and the standard √ deviation is  = .

SECTION 27.3

PROBLEMS

Binomial Distribution 1. In each of (a) through (f), a number N of trials is performed, each with probability p of success. Use a binomial distribution to compute the probability of x successes for the given values. (a) N = 8 p = 043 x = 2 (b) N = 4 p = 07 x = 3 (c) N = 6 p = 05 x = 3 (d) N = 10 p = 06 2 ≤ x ≤ 5 (e) N = 8 p = 04 x ≥ 7 (f) N = 10 p = 058 2 ≤ x ≤ 4 (g) N = 10, p = 035, x = 3 or x = 7 (h) N = 7, p = 024, x = 1 or x = 3 or x = 5 2. A professor gives a test consisting of six multiple choice questions, each offering four answers from which to choose. A student who knows less than nothing takes the test by the method of random choice, simply guessing on each question.

(a) What is the probability that the student gets exactly one question right? (b) What is the probability of getting exactly two right? (c) What is the probability of getting four or five right? (d) What is the probability of getting all six questions right? 3. In Problem 2, change the number of choices for each question to three. Now answer questions (a) through (d). 4. In Problem 2, keep the number of choices given with each question at four, but suppose now there are ten questions. Answer the questions posed in Problem 2, and also determine the probability of getting five, six, or seven questions correct. 5. An auto manufacturer purchases starting mechanisms from an independent contractor. Every week, a quality control officer pulls one hundred starters out of the shipment from the contractor and inspects each. A starter is either defective or not, with no range in

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve between. If two or fewer are defective, the total shipment of the week is accepted. Otherwise it is returned to the contractor. Further, it is known that overall the starter company has a 32% rate of production of defective products. What is the probability that the next shipment will be accepted? 6. In a clinical trial, it was found that 7% of those in the control group (who did not receive the drug) experienced a particular side effect commonly associated with this drug. Suppose the same 7% rate was found in those taking the drug in the trial. Find the probability that among the 1000 subjects taking the drug exactly 92 exhibited the side effect.

Poisson Distribution 7. For  = 13, plot points of the Poisson distribution corresponding to x = 0 1 2     8. Draw a smooth curve through these points. This will gives some sense of the fact that a Poisson distribution has an approximately bell-shaped curve although here x ≥ 0 and the curve is not an exact bell curve. 8. Assuming a Poisson distribution, determine P
x for each of the following. (a)  = 09 x = 6 (b)  = 085 x = 10 (c)  = 092 x = 4 (d)  = 087 x = 8 (e)  = 094 1 ≤ x ≤ 5

9. In a splatter test conducted over a full day, a car windshield is subdivided into 320 pieces, and it is found that 295 were hit by some material kicked up from the road by other cars. Assuming a Poisson distribution, determine (a) the probability that some piece is hit three times (b) the probability that some piece is hit from two to five times, inclusively 10. A new computer chip has one million microdots etched on its surface. In the manufacturing process, 997 850 dots are imprinted with information. Assuming a Poisson distribution, determine the probability that the process will incorrectly etch exactly three information packets on the same dot. 11. A golf driving range subdivides the area between 100 and 300 yards from the tee into 500 regions. In one four hour period, 476 balls were driven into this area. Assuming a Poisson distribution, compute the probability that any region was hit 3 times. What is the probability that a region was hit from two to six times, inclusive? 12. A gaming machine randomly selects a positive integer each time fifty cents is inserted. A trained cocker spaniel always bets on 15. The probability of 15 coming up is 1/24. The spaniel plays the game 100 times. Assuming a Poisson distribution, find the probability

(f)  = 064 3 ≤ x ≤ 5

(a) of winning 7 times

(g)  = 075 x = 3 or x = 8

(b) of winning 3 times

(h)  = 097 x = 1 or x = 3 or x = 10

(c) of not winning even once

27.4

1159

A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve Here is a simple but instructive experiment. Flip a coin twelve times and record the number of heads. If we perform this experiment once, some number of heads from 0 to 12, inclusively, will come up. We would be surprised to see no heads or twelve heads, but this is not impossible. An outcome of, say, five, six or seven heads would not cause much reaction. Repeat the experiment (12 flips) a few times, recording the number of heads each time. On any given set of 12 flips, we can get any number of heads from 0 to 12, inclusive. Now suppose we conduct N sets of 12 flips. Record the number of heads on each set of 12 flips. Is there any pattern that becomes apparent in carrying out large numbers of repetitions of 12 flips? Since in any repetition the number of heads can be any integer from 0 through 12, we might at first guess that there is not. However, a perhaps surprising regularity begins to occur if we look not just at outcomes, but at frequencies of outcomes. Instead of making a list of

CHAPTER 27

Statistics

the number of heads in each repetition, count, over all the repetitions, the frequency of the occurrence of any given number of heads. For example, suppose we run 50 repetitions of 12 flips. If five heads occur in exactly 15 of these repetitions, we say that 5 has a frequency of 15. If eight heads occur in exactly 12 of the repetitions, then 8 has a frequency of 12. In this way, we can look at the frequency of each outcome over all the repetitions. To illustrate, below are the results of N = 21 repetitions of the 12 coin flips, recording the number of heads on each set. (If you carry out your own 21 sets, you may get different outcomes). The first set of twelve flips had four heads, the seventh set eight heads, the thirteenth had six heads, and so on. Exactly three of the sets had four heads, so four heads has frequency 3. Similarly, 8 heads occur six times in row two of Table 27.1, so 8 heads has a frequency of 6, and 9 heads has a frequency of 3, and so on. Table 27.2 is a frequency table, summarizing the results of Table 27.1 in terms of frequencies: It is convenient to display this frequency table in a graph of number of heads (horizontal) against frequency (vertical). The resulting bar graph is shown in Figure 27.4. This bar graph, based on only 21 repetitions, does not suggest any pattern. However, Figures 27.5 through 27.9 show, in turn, frequency graphs for 86 254 469 920 and, finally,

TA B L E 2 7 . 1 Set of Twelve Flips Number of Heads

1 2 3 4 5 6 7 8 9 10 11 12 13 14 15 16 17 18 19 20 21 4 6 9 8 5 6 8 6 8

9

8

6

6

8

8

6

3

9

4

4

6

TA B L E 2 7 . 2 Number of Heads Frequency of this Number in 21 Trials

0

1

2

3

4

5

6

7

8

9

10

11

12

0

0

0

1

3

1

7

0

6

3

0

0

0

4

5 6 7 Number of heads

11

12

8 7 6 Frequency

1160

5 4 3 2 1 0 1

FIGURE 27.4

2

3

12 coin flips, 21 repetitions.

8

9

10

Frequency

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1161

20 18 16 14 12 10 8 6 4 2 0 1

2

3

4

5 6 7 Number of heads

8

9

10

11

12

FIGURE 27.5 12 coin flips, 86 repetitions.

60

Frequency

50 40 30 20 10 0 1

2

3

4

5 6 7 8 Number of heads

9

10

11

12

FIGURE 27.6 12 coin flips, 254 repetitions.

120

Frequency

100 80 60 40 20 0 1

2

3

4

5

6 7 8 Number of heads

FIGURE 27.7 12 coin flips, 469 repetitions.

9

10

11

12

CHAPTER 27

Statistics

250

Frequency

200 150 100 50 0 1

2

3

4

5

6

7

8

9

10

11

12

Number of heads FIGURE 27.8

12 coin flips, 920 repetitions.

350 300 250 Frequency

1162

200 150 100 50 0 1

2

3

4

5

6

7

8

9

10

11

12

Number of heads FIGURE 27.9

12 coin flips, 1414 repetitions.

1414 repetitions of the 12 flips. Figure 27.10 shows a summary of the bar charts of Figures 27.5 through 27.9, except that, instead of bars, we have drawn a curve through the tops of the bars. The bottom curve is for 86 repetitions of the 12 flips, the next for 254 repetitions, and so on through 469, 920, and—the highest curve—1414 repetitions. Notice that, as the number of repetitions increases, the frequency graphs in figure 27.10 approach a bell curve. We will give a more careful characterization of bell curves shortly, but for now, this is a curve having the general appearance of the graph of Figure 27.11. The fact that the frequency charts approach a bell-shaped curve as N increases is characterized by saying that the data of these trials are normally distributed. Normal distribution means that there is a bell curve which can be fit (approximately) to the data. We will see that such a bell curve contains a great deal of information about probabilities of the experiment. Part of our task will be to see how to extract this information. There are many other normally distributed data sets whose graphs approximate bell curves. For example, Table 27.3 gives score ranges on a standardized test that is administered nationally, and the percentage of the students taking the test who scored in each range that year.

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1163

350 300 250 200 150 100 50 0 1

2

3

4

5

6

7

8

9

10

11

12

FIGURE 27.10 Line graphs of frequencies of heads for 86, 254, 469, 920 and 1414 repetitions.

0.2

0.15

0.1

0.05

0 –4

0

4

8

12

FIGURE 27.11 Typical bell curve.

This data is graphed as a bar chart in Figure 27.12 (a). If a curve is drawn through the tops of the bars, as in Figure 27.12 (b), an approximately bell-shaped curve results. We have now seen two examples of curves we have referred to as "bell curves". Here is a definition of this term.

CHAPTER 27

Statistics

TA B L E 2 7 . 3 Range 200–240 250–290 300–340 350–390 400–440 450–490 500–540 550–590 600–640 650–690 700–740 750–800

Percent of Students in this Range 09 18 46 89 135 176 175 143 106 64 25 14

20 Percent of students in range

1164

18 16 14 12 10 8 6 4 2 0 0 290 340 390 440 490 540 590 640 690 740 800 0– 00– 50– 00– 50– 00– 50– 00– 50– 00– 50– 25 3 3 4 4 5 5 6 6 7 7

24

0–

20

Score ranges FIGURE 27.12(a)

Standardized test scores.

20 18 16 14 12 10 8 6 4 2 0 0

2

FIGURE 27.12(b)

4

6

8

10

Standardized test scores bell curve.

12

14

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

DEFINITION 27.8

1165

Bell Curve

A bell curve is the graph of an exponential function y= √

1

e−
x−

2 /2 2

2



in which  is any positive number and  is any number.

Figure 27.13 shows graphs of three bell curves drawn on the same axes for comparison. All curves have  = 6, and the choices for  are 35, 4 and 45. The following are some important properties shared by every bell curve. 1. The curve lies entirely above the horizontal axis and and has a single maximum point at x = . 2. The curve is symmetric about the vertical line x = . 3. The curve is asymptotic to the horizontal axis as x → and as x → − . 4. The curve has exactly two points of inflection. 5. The area bounded by the bell curve and the x-axis is 1: 1
−
x−2 /2 2 e dx = 1 √ 2 − We will see that  is actually a standard deviation, and  is a median. 0.7

σ = 3.5

0.6

σ=4

0.5

σ = 4.5

0.4

0.3

0.2

0.1

0 0

5

10

FIGURE 27.13 Bell curves corresponding to  = 6 and  equal to 35, 4 and 45.

15

1166

CHAPTER 27

Statistics

It is Property 5 which makes possible the following connection between a bell curve and probability.

THEOREM 27.2

Suppose the bell curve y = Then,

2 2 √ 1 e−
x− /2 2

is the graph of a continuous random variable X.

1. If a < b, and P
a ≤ x ≤ b denotes the probability that an observed data value X
x falls within a b, then 1
b −
x−2 /2 2 P
a ≤ x ≤ b = √ e dx 2 a (See the shaded portion of Figure 27.14.) 2. If P
x ≥ a denotes the probability that an observed data value X
x is at least as large as a, then P
x ≥ a = √

1




2



e−
x−

2 /2 2

dx

a

(See the shaded portion of Figure 27.15.) 3. If P
x ≤ b denotes the probability that an observed data value X
x is no larger than b, then P
x ≤ b = √

1 2




b

−

e−
x−

2 /2 2

(See the shaded part of Figure 27.16.)

a FIGURE 27.14

b

Shaded area equals P
a ≤ x ≤ b.

dx

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1167

a FIGURE 27.15 Shaded area equals P
x ≥ a.

b FIGURE 27.16 Shaded area equals P
x ≤ b.

A few years ago these integrals would be approximated using tables constructed for this purpose. Now, however, the integrals are quickly approximated by mathematics computation software, such as MAPLE or MATHEMATICA. Now we must make the significant leap from a continuous random variable to a discrete random variable, since the latter are what are often generated by experiments or procedures that we actually carry out. We have seen two examples (coin flips and/standardized test scores) in which the data approximated a bell curve. Such data, which are often values of a discrete random variable, are said to be normally distributed. We will now show how the mean and standard deviation

1168

CHAPTER 27

Statistics

of normally distributed data can be used to associate the data with a particular bell curve (by specifying the constants  and ), and then how conclusions about the probability of certain outcomes related to the data can be drawn from the curve, using the theorem. Suppose we carry out some set of trials (such as flipping a coin, or inspecting items for defects, or determining whether a drug has beneficial results), and each trial has exactly two possible outcomes, which we will call A and B, or success and failure, or some other bimodal designation. Suppose also the outcomes of separate trials are independent of each other. We therefore have a binomial random variable which we will call X, where X
x = number of times A occurs in x. Let the probability that A occurs in any trial be p. (For an honest coin flip, p = 1/2). By complementarity the probability that B occurs is q = 1 − p. The frequency of A after N trials, as N increases, approaches a normal distribution given by a bell curve. The mean of the frequencies of any set of N trials given by this bell curve is  = Np and the standard deviation of this data is =



Np
1 − p

The bell curve, or normal curve, associated with these trials is the graph of 1 2 2 y= √ e−
x− /2 2

(27.3)

Now suppose we pick an interval of values on the x-axis, and we want the probability that X
x falls within this interval (that is, that A occurs between a and b times, inclusive, in x). Here we must be careful. The graph of the continuous bell curve from equation (27.3) is only approximated by the graph of the discrete random variable. We must therefore make an adjustment in relating area under this curve to a probability. This is done by what is called a continuity adjustment in the limits of integration, effected by adding 1/2 to the upper limit of integration, and subtracting 1/2 from the lower limit. This is described in the following theorem. THEOREM 27.3

1. Let a ≤ x ≤ b. Assume that the random variable X has a binomial distribution. Then P
a ≤ x ≤ b, denote the probability that an observed data value X
x falls within a b, is given by
b+1/2 1 2 2 P
a ≤ x ≤ b ≈ e−
x− /2 dx √ a−1/2 2 2. P
x ≥ a = probability that the observed value is at least as large as a
1 2 2 ≈ e−
x− /2 dx √ a−1/2 2 3. P
x ≤ b = probability that the observed value is no larger than b
b+1/2 1 2 2 ≈ e−
x− /2 dx √ − 2

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1169

Keep in mind that this continuity adjustment in the limits of integration applies only to the bell curve approximation of a binomial distribution. The reason these conclusions are stated as approximations is that we are approximating the discrete random variable X by a continuous bell curve (chosen with  and  determined by X), and also the limits of integration include a continuity adjustment which is itself an approximation. We will illustrate how the theorem is used to extract numerical information about events in an experiment.

EXAMPLE 27.16

Flip a coin 12 times. Previously we performed many repetitions of the 12 flips, charting the frequencies of the occurrences of heads. Here N = 12, and the probability of obtaining a head on any one flip of the coin is p = 1/2, assuming an honest coin. The only other outcome of any coin toss is a tail, which has probability q = 1 − 1/2 = 1/2. The mean of the (approximately) normally distributed data is 1  = Np =
12 = 6 2 not a surprising result. Ideally, if we flip a coin 12 times, we expect on average to obtain 6 heads, although in a specific run of 12 flips we might get any number of heads from 0 through 12. The standard deviation is   = Npq =

0 12

   √ 1 1 = 3 2 2

The bell curve for this experiment is y= √

1

e−
x−

2 /2 2

2 1 1 2 2 = √ √ e−
x−6 /6 = √ e−
x−6 /6  2 3 6

A graph of this curve is shown in Figure 27.17. The maximum occurs at x = 6, and this line is also the axis of symmetry. Suppose we want to compute the probability that the number of heads that comes up falls between, say, 4 and 8. Taking care to use the continuity adjustment, this probability is 1
85 −
x−62 /6 P
4 ≤ x ≤ 8 ≈ √ e dx ≈ 085109 6 35 The probability is approximately 085 that a head will come up 4, 5, 6, 7 or 8 times. Figure 27.18 shows this probability as an area under the approximating normal curve.

1170

CHAPTER 27

Statistics

0.20

0.15

0.10

0.05

0 0 FIGURE 27.17

2

4

6

8

10

12

8

10

12

1 2 Bell curve y = √ e−
x−6 /6 . 6

0.20

0.15

0.10

0.05

0 0 FIGURE 27.18

2

4

6

P
4 ≤ x ≤ 8, with the continuity adjustment.

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1171

In this example we can compute P
4 ≤ x ≤ 8 exactly using a binomial distribution for comparison: P
4 + P
5 + P
6 + P
7 + P
8    4  8    5  7    6  6 1 1 1 1 1 1 12 12 12 = + + 4 5 6 2 2 2 2 2 2    7  5    8  4 1 1 1 1 12 12 + + 7 8 2 2 2 2           1 12 12 12 12 12 + + + + = 12 4 5 6 7 8 2 =

1
495 + 792 + 924 + 792 + 495 = 0854 212

differing from the approximation by about three thousandths. If we left out the continuity adjustment, we would compute P
4 ≤ x ≤ 8 as 1
8 −
x−62 /6 e dx √ 6 4 which is approximately 0751 79. We know that this is incorrect because we have computed P
4 ≤ x ≤ 8 exactly in this case using a binomial distribution. What is the probability that the number of heads will be 10 or more? This is 1
2 P
x ≥ 10 ≈ √ e−
x−6 /6 dx ≈ 002165 10−1/2 6 As we might expect, getting 10 or more heads in 12 flips is very unlikely. This probability is represented by the shaded area in Figure 27.19. What is the probability of getting from 0 to 6 heads, inclusive? This is the area of the shaded region in Figure 27.20, and is given by 1
6+1/2 −
x−62 /6 e dx ≈ 061359 P
x ≤ 6 ≈ √ 6 − It is a common error to think that P
x ≤ 6 should be 1/2. However, there are seven ways the number of heads can be 6 or less, but only six ways it can be more than 6. This probability also dramatizes the need for the continuity adjustment. The area under this bell curve, for x ≤ 6, is indeed half the total area bounded by the curve, or 1/2. However, the probability P
x ≤ 6 is not 1/2. The apparent discrepancy is due to the fact that this bell curve is only an approximation to the data of the experiment. Finally, what is the probability that there will be exactly 7 heads? This probability is 1
7+1/2 −
x−62 /6 P
x = 7 = √ e dx ≈ 0193 6 7−1/2 Again, this is a probability we can compute without a bell curve as  P
7 =

12 7

  7  5 1 1 = 0193 · · ·  2 2

1172

CHAPTER 27

Statistics

0.20

0.15

0.10

0.05

0 0 FIGURE 27.19

2

4

6

8

10

12

10

12

P
x ≥ 10, with the continuity adjustment.

0.20

0.15

0.10

0.05

0 0 FIGURE 27.20

2

4

6

8

P
x ≤ 6, with the continuity adjustment.

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1173

EXAMPLE 27.17

Suppose we roll a die 50 times. We are interested in the frequency of having 4 come up. Think of each roll or trial as having two outcomes of interest, either 4 or not 4. Since there are six faces on the die, the probability of getting a 4 on one roll is p = 1/6, while the probability of getting not 4 is 5/6. Since N = 50 for fifty trials, the mean of the normally distributed data is 1 25  =
50 =  6 3 approximately 833. The standard deviation is 0     5 250 5 √ 1  = 50 = = 10 6 6 36 6 approximately 264. The bell curve associated with these trials is the graph of y= √

1 2

e−
x−

2 /2 2

1 2 = √  √ e−9
x−25/3 /125 2 56 10

6 2 e−9
x−25/3 /125  = √ 5 20 Suppose, for example, we want the probability that, on fifty rolls of the die, the number of times 4 comes up is between 2 and 15, inclusively. This probability is (approximately) the

0.14

0.12

0.10

0.08

0.06

0.04

0.02

0 0

4

8

FIGURE 27.21 P
2 ≤ x ≤ 15, using the bell curve

6 2 e−9
x−25 /125 . y= √ 5 20

12

16

1174

CHAPTER 27

Statistics

shaded area under the bell curve shown in Figure 27.21. This area appears to be nearly all of the area under the curve, so we expect P
2 ≤ x ≤ 15 to be close to 1. Compute
15+1/2 6 2 P
2 ≤ x ≤ 15 ≈ √ e−9
x−25/3 /125 dx ≈ 099198 5 20 2−1/2 Suppose we want the probability that the number of 4’s that come up in fifty rolls is 47 or 48. The area between the bell curve and the horizontal axis for x between 47 and 48 is nearly 0. In fact, it is so small that it doesn’t show up on the graph (in the scale of the drawing). This means that, although 47 or 48 4s are not impossible in fifty rolls, it is extremely unlikely that either of these number of 4s will come up. To check this intuition, compute
48+1/2 6 2 P
47 ≤ x ≤ 48 ≈ √ e−9
x−25/3 /125 dx 5 20 47−1/2 ≈ 77325
10−48 

EXAMPLE 27.18

Suppose an airline finds (by analyzing past flight information) that its planes fly an average of 1500 miles per month per passenger, with a standard deviation of 90 miles. Suppose we want the probability that a typical passenger flies between 1400 and 1600 miles per month. Under the assumption that the data is normally distributed, we have  = 1500 and  = 90, and these determine the appropriate bell curve for this study: y= √

1 2

e−
x−

2 /2 2

=

1 √

90 2

e−
x−1500

2 /16200



The probability that a typical passenger flies between 1400 and 1600 miles per month is approximately P
1400 ≤ x ≤ 1600 ≈

1 √




90 2

1600+1/2

e−
x−1500

2 /16200

dx

1400−1/2

≈ 0735 86

27.4.1 The Standard Bell Curve The bell curve corresponding to  = 0 and  = 1 is the graph of 1 2 y = √ e−z /2  2

(27.4)

The graph is shown in Figure 27.22, with the horizontal axis as the line of symmetry. This curve is called the standard or normal bell curve. It is routine to use z as the variable in this context, and statisticians often refer to values on the horizontal axis as z-scores. One reason for having a standard bell curve is that it provides a reference point from which to compile certain information in tables. To illustrate one case where this is useful, suppose we have in mind a particular area under the standardized bell curve. For example, we may have a particular target probability that we are aiming for, and we want to know a value z = z0 so

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1175

0.4

0.3

0.2

0.1

0 –3

–2

–1

0

1 1

FIGURE 27.22 Standard bell curve y = √

2

2 2 /2

e−z

3

.

that the area under the curve, to the left of z0 , gives that probability. If the probability we are aiming for is k, then we need to solve 1
z0 −z2 /2 e dz = k √ 2 − for z0 , a highly nontrivial task. Table 27.4 gives some commonly used probabilities, and the value of z0 such that that the area to the left of z0 under the standard bell curve gives this probability. For example, 1
1645 −z2 /2 e dz = 095002 √ 2 −

TA B L E 2 7 . 4 Area to the Left of z0 099 095 090 085 080 075

z0 233 1645 129 104 084 067

1176

CHAPTER 27

Statistics 0.4

0.3

0.2

0.1

0 –3

–2

–1

0

1

2

3

z0 = 1.29 FIGURE 27.23 Area under the standard bell curve to the left of z0 = 129 is approximately 0.90.

and 1
129 −z2 /2 e dz = 090147 √ 2 − The last integral is illustrated in Figure 27.23.

27.4.2

The 68, 95, 997 Rule

Suppose we have some experiment that involves a fixed number of trials. For example, in flipping a coin fifty times, the experiment consists of the fifty flips, each of which is a trial. Let the probability of the result of a trial that we are interested in(such as heads for a coin toss) be p. Compute the mean  = Np and standard deviation  = Np
1 − p. Define the random variable X
x = number of occurrences of the trial result of interest in outcome x. Then the bell curve associated with X is 1 2 2 y= √ e−
x− /2  2

(27.5)

27.4 A Coin Tossing Experiment, Normally Distributed Data, and the Bell Curve

1177

It is routine to check (solve y


x = 0) that this curve has two points of inflection, namely x =  +  and x =  − . The points of inflection are therefore one standard deviation to the left and right of the mean, hence to the left and right of the bell curve’s axis of symmetry. Moving one, two or three standard deviations to the left or right of the mean is particularly significant, according to the following rule. The 68 95 997 Rule The area under the bell curve between  −  and  +  is approximately 68% of the total area. The area under bell curve between  − 2 and  + 2 is approximately 95% of the total area. The area under bell curve between  − 3 and  + 3 is approximately 997% of the total area. In terms of probabilities: The probability that the outcome falls between  −  and  +  is approximately 068. The probability that the outcome falls between  − 2 and  + 2 is approximately 095. The probability that the outcome falls between  − 3 and  + 3 is approximately 0997 The probability that an outcome falls to the left of  − 3 or to the right of  + 3 is 1 − 0997 = 0003. Hence it is nearly certain that any outcome of the experiment will fall within three standard deviations of the mean.

EXAMPLE 27.19

Suppose we have a dishonest coin, and the probability of a head on any toss is p = 08. Flip the coin 1000 times. Compute  = 1000
08 = 800 and   = 1000
08
02 = 1265 Now  −  = 78735 and  +  = 81265, so in the thousand flips, there is a probability of about 68 of seeing between 787 and 813 heads. There is a probability of about 95 of seeing between 774 and 826 heads, and a probability of about 997 of seeing between 762 and 838 heads.

SECTION 27.4

PROBLEMS

1. Roll a die 90 times. We are interested in the outcome that a 6 comes up. (a) Determine  and  for this experiment. (b) Write the equation of the corresponding bell curve and graph the curve. Determine the probability that a 6 comes up:

(f) fewer than 20 times (g) exactly 45 times 2. A coin is flipped 100 times. (a) Determine  and  for this experiment.

(c) between 30 and 60 times

(b) Write the equation of the appropriate bell curve and graph the curve. Determine the probability that a head comes up:

(d) between 2 and 80 times

(c) between 20 and 40 times

(e) at least 35 times

(d) between 10 and 50 times

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(e) between 45 and 55 times (f) at least 30 times (g) fewer than 60 times (h) exactly 55 times 3. In a certain kingdom, babies are born with a probability of 048 for boys and 052 for girls. In a particular year, it is expected that 350 children will be born. We are interested in the outcome that a newborn child was a girl. (a) Determine  and  for this experiment. (b) Write the equation of the bell curve and graph the curve. Determine the probability that: (c) at least 220 of the babies will be girls (d) at least 150 of the babies will be girls (e) the number of girl babies will fall between 120 and 250 (f) the number of girl babies will be exactly 180 4. It is known that, when a stratigather is produced by a certain company, it may be either left handed or right handed. The probability that a particular stratigather will be right handed is 058, and that it will be left handed, 042. A thousand stratigathers are picked at random. Compute the probability (a) That at least 400 of these are right handed. (b) That between 400 and 600 are right handed. (c) That no more than 450 are right handed. (d) That exactly 520 are right handed. 5. A bus company finds that its buses drive an average of 700 miles per month for each passenger with a standard deviation of 65 miles. Assume that the data consisting of the number of miles flown per passenger per month

27.5

is normally distributed. Determine the probability that a typical passenger rides between 250 and 600 miles per month. What is the probability that the passenger rides between 600 and 900 miles per month? 6. An auto rental firm finds that its cars travel an average of 940 miles for each rental, with a standard deviation of 76 miles. Find the probability that a customer who rents the car drives it between 750 and 1000 miles. What is the probability that the car is driven between 300 and 500 miles? 7. In a recent year in the American League, hitters averaged 247, with a standard deviation of 021. Determine the probability that a hitter averaged between 245 and 270. What is the probability that a hitter averaged over 260? What is the probability that a batter averages 300 or more? 8. A coin has a probability of 042 of coming up heads if it is flipped. The coin is flipped 2000 times. (a) Determine a and b so that there is a probability of about 068 of seeing between a and b heads. (b) Determine a and b so that there is a probability of about 095 of seeing between a and b heads. (c) Determine a and b so that there is a probability of about 099 of seeing between a and b heads. 9. An honest die is rolled, and we are interested in the outcome that the number that comes up is even. A total of 550 tosses are made. (a) Determine a and b so that there is a probability of about 068 of seeing between a and b even tosses. (b) Determine a and b so that there is a probability of about 095 of seeing between a and b even tosses. (c) Determine a and b so that there is a probability of about 099 of seeing between a and b even tosses.

Sampling Distributions and the Central Limit Theorem We now have some facts and information about data sets and how certain statistics can be extracted from them. We now move toward addressing a fundamental problem of statistics. In statistical studies, the term population refers to the set of all the objects under study. For example, if we are studying airplane parts coming off a factory production line in a certain time period, then the population would consist of all the parts produced in that time. Suppose we want to measure the success of some process (such as manufacturing) by estimating how many defective products are produced. An obvious strategy is to sample the product to try to infer from the sample a measure of how many products are defective in the entire population. How can we do this in a reasonable and efficient manner? Clearly the answer lies somewhere between the extremes of inspecting every item (probably costly and inefficient), and inspecting only none or one (of questionable value). How do we decide how many items

27.5 Sampling Distributions and the Central Limit Theorem

1179

to sample, and what confidence can we have in conclusions drawn from samples? In particular, can we tie the number of samples to a predetermined level of confidence, so that we can set the sample size to achieve this confidence level? To approach this important issue, let us agree that, in taking a sample from a collection of objects, we will sample with replacement. This means that, if we pick an object and make some measurement on it, then we replace this object back in the population before picking another object to test. This way each sample is chosen from the entire population. One reason for doing this is that usually very large numbers of objects are involved, so sampling with or without replacement makes little numerical difference. In addition, sampling with replacement ensures that the selection of each sample object is independent from the selection of any other, since all objects are chosen from the entire set of objects. This is a considerable simplification. We also will agree that sample size is fixed at some number n throughout the discussion. Thus, all samples in a given study are taken with replacement, and have the same sample size. When we take samples involving numerical quantities, each sample will have a mean (mean of the numbers in the sample). The sampling distribution of the mean is the probability distribution of the means of the samples, assuming that all the samples have the same size. We also may compute other statistics from the samples, such as medians of the samples and standard deviations of the samples, and average these.

EXAMPLE 27.20

A contractor looks at the number of buildings his company has completed over the past five years, and the numbers are 1, 4, 6, 8 and 9. These numbers constitute the population under consideration. Samples of size two are taken, with replacement. Of course in reality, we would not take samples of two out of five, but this example has the sole purpose of illustrating the discussion. With replacement, there are 52 = 25 samples of two of the numbers. The samples are listed in Table 27.5, together with each sample mean and sample standard deviation. If we average all fifty numbers contained in the samples, we find that this mean is 56. Column 2 of the table contains the mean of each sample, which in this case is an average of just two numbers for each sample. If the means x for each sample x are averaged (that is, average the numbers in column two), we find that this sample mean is also 56. This is not a coincidence. The mean of all the numbers in the samples is the same as the average of all the sample means. We say that this sample statistic (mean of the samples) targets the population parameter (mean of all the items appearing in the samples). The sampling distribution of the mean is the probability distribution of the sample means (which are in Column 2). There are 25 samples, and, for example, the number 5 occurs four times as an average of samples, so P
5 = 4/25. Similarly, the number 35 occurs twice as an average of samples, so P
35 = 2/25. The entire sample distribution is: P
1 =1/25 P
25 = P
35 = P
45 = 2/25 P
5 = 4/25 P
4 =1/25 P
6 = 3/25 P
65 = P
7 = P
75 = 2/25 P
8 = 1/25 P
8 = 2/25 P
9 = 1/25 As required, these probabilities sum to 1. If we compute the mean of the sample standard deviations (Column 3), we get 22368. However, if we compute the standard deviation of the fifty items in the samples, we get 28705. The standard deviation of the samples is not targeted by the mean of the sample standard deviations.

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TA B L E 2 7 . 5 Sample x 1,1 1,4 1,6 1,8 1,9 4,1 4,4 4,6 4,8 4,9 6,1 6,4 6,6 6,8 6,9 8,1 8,4 8,6 8,8 8,9 9,1 9,4 9,6 9,8 9,9

Sample Mean x 1 25 35 45 5 25 4 5 6 65 35 5 6 7 75 45 6 7 8 85 5 65 75 85 9

Sample Standard Deviation 0 21213 35355 43012 56569 21213 0 14142 28284 35355 35355 14142 0 14142 21213 49497 28284 14142 0 07071 56569 35355 21213 07071 0

In making inferences about a population by observing samples, the sample distribution of a proportion often plays a role. Suppose a sample of n objects from the population has k of one type of object. Then the sample proportion of that sample for that item is k/n. The sampling distribution of the population, with respect to that item, is the probability distribution of the sample proportions for that item.

EXAMPLE 27.21

Suppose, in trying to determine quality of manufacture of a product, ten samples of 20 items each are taken off the line (with replacement). Table 27.6 shows the data and the sample ratios (number defective divided by sample size): The sample ratio for each sample (with regard to defective items) is the number of defective items in the sample, divided by n = 20. Since 1/20 occurs four times out of ten as a sample ratio, P
1/20 = 4/10 = 04. In this way we compute the sampling distribution of this proportion: P
1/20 = 04 P
0 = 04 P
1/10 = 02

We now have the vocabulary to state the important central limit theorem, which relates a distribution of the original population itself to a distribution of the sample means.

27.5 Sampling Distributions and the Central Limit Theorem

1181

TA B L E 2 7 . 6 Sample 1 2 3 4 5 6 7 8 9 10

THEOREM 27.4

Number Defective 1 0 2 0 0 1 2 1 1 0

Sample Ratio 1/20 0 1/10 0 0 1/20 1/10 1/20 1/20 0

Central Limit Theorem

Suppose a random variable X has mean  and standard deviation . Then the distribution of sample means will approach a normal distribution as the sample size increases, regardless of whether or not the original population distribution is normally distributed. Further, 1. The mean of all the sample means equals . √ 2. The standard deviation of all sample means is approximately / n, where n is the sample size. This approximation can be improved by increasing n. The first conclusion can be phrased: the mean of the samples equals the mean of the population. The second conclusion justifies the sense in which we can approximate the standard deviation of the total population using the standard deviation of the sample means. The central limit theorem justifies the approximation of a distribution by a normal distribution. Even when the original random variable does not have a normal distribution, the sample means will approach a normal distribution. As a rule of thumb (that is, true except for “pathological” cases), we can safely approximate the sample means distribution by a normal distribution if n > 30. The following example illustrates the idea of the central limit theorem, and emphasizes the crucial difference that must be understood between the distribution of the original random variable, and the distribution of the sample means.

EXAMPLE 27.22

Imagine that forty samples of six numbers each have been drawn, with replacement, from the integers 0 through 9, inclusive. Table 27.7 shows the samples and the sample means. The sample column has forty entries of six numbers each, hence a total of 240 numbers (the population), each an integer from 0 through 9, inclusively. The random variable X is defined by X
k = number of occurrences of k in this set of 240 numbers. Figure 27.24 shows a graph of these values of X
k. Clearly this is not a normal distribution, nor should we expect it to be. Figure 27.25 shows a graph of the distribution of the sample means. This distribution is not normal, because n = 6 is too small. However, if n were chosen larger, the sample means would more closely approximately a true normal distribution.

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TA B L E 2 7 . 7 Sample 289186 957888 231332 767099 686505 859901 063159 501547 855909 161078 809973 290800 925079 732111 795556 145776 276535 605812 633966 747358

Sample Mean 566 75 233 633 5 533 4 366 6 383 6 316 533 25 616 5 466 366 550 560

Sample 188635 405396 148235 332674 962669 445186 637269 521802 040848 358298 568735 630761 607495 188965 201037 749018 180629 291554 958087 799365

Sample Mean 516 45 383 416 633 466 550 3 4 583 560 383 516 616 216 483 433 433 616 65

35 30 25 20 15 10 5 0 0 FIGURE 27.24

2

4

6

8

10

Graph of X(k).

It is routine to compute the mean of the population, obtaining =

387  80

A similar routine calculation shows that the mean of the sample means is also 387/80, as expected. The standard deviation of the population is found to be =

√ 492 159  240

27.5 Sampling Distributions and the Central Limit Theorem

1183

14 12 10 8 6 4 2 0 0

2

4

6

8

FIGURE 27.25 Distribution of the sample means.

which is approximately 29231. The standard deviation of the sample means is √ 86 919  means = 240 √ √ or approximately 12284. If we compute / n, which is / 6 for this example, we get 11933. According to the central limit theorem, this should approximate means , and indeed this number is within 0035 of means . Even with a small n, this approximation is good to about three hundredths in this example.

EXAMPLE 27.23

The headmaster of an exclusive private high school is interested in the probability that all twenty two of his seniors will score at least 670 out of 900 on a standardized test used by colleges in making admission decisions about applicants. The national average on this test is 640, with a standard deviation of 120. Think of taking samples of size n = 22 from the national pool of students taking this test, and use the approximately distributed sample means to draw conclusions about this sample. For the entire population,  = 640 and  = 120. For the sample means, means =  = 640 and 120 means = √  22 The bell curve for the sample means is the graph of √ 2 1 2 y =√ e−
x−640 /2
120/ 22 √ 2
120/ 22 √ 11 −11
x−6402 /14400 =  √ e 120 

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0.012

0.008

0.004

0 600

620

FIGURE 27.26

x ≥ 670.

640

660

680

Area under the bell curve y =

700

720

740

√ 2 11 √ e−11
x−640 /14400 120 

for

We want the area under this bell curve for x ≥ 670. This area is shown in Figure 27.26, and is approximately equal to √ 11
900 −11
x−6402 /14400 e dx √ 120  670 or 0120 48. Perhaps on average, the headmaster can expect the entire class to score this high once in ten years, although it also could happen two years in a row, or never in this headmaster’s tenure. What is the probability that all twenty-two students of the students score between 640 and 700? This is √ 11
700 −11
x−6402 /14400 e dx √ 120  640 or 0490 49, very close to one half. What is the probability that ten of the students score better than 660? Now use n = 10 in forming the samples, so the bell curve is √ 5 2 y= √ e−5
x−640 /14400  120  and the probability we want is

which is 0299 08, about 3/10.

√ 5
900 −5
x−6402 /14400 e dx √ 120  660

27.6 Confidence Intervals and Estimating Population Proportion

SECTION 27.5

PROBLEMS

1. Thirty samples are taken of five digit numbers. Each of these numbers has been drawn, with replacement, from the integers 3 through 9, inclusive. The samples are given in the following table: Sample

Sample

Sample

Sample

Sample

74974 76693 43856 85994 36496 48659

78365 57733 86934 43865 59378 43958

94489 58438 96649 67339 47937 48896

84937 59786 88533 59643 75488 69537

33933 58393 48576 58834 84938 78953

(a) Compute the sample means, the mean of the sample means, and the mean of the population. (b) Compute the standard deviation of the population and the standard deviation √ of the sample means, and verify that means ≈ / n for the appropriate choice of n. 2. Thirty samples are taken of seven digit numbers. Each of these digits has been drawn, with replacement, from the integers 0 through 9, inclusive. The samples are given in the following table: Sample 9388204 8902817 7921082 8019283 5501829 4039112

Sample 7939720 3022850 5910293 4911029 3371932 9571120

Sample 5773902 5492018 9017839 3019402 0002910 9912118

Sample 1912038 4491002 4491029 5616209 3611390 8401938

Sample 1749823 4910292 4019283 9201898 4092179 7109032

(a) Compute the sample means, the mean of the sample means, and the mean of the population. (b) Compute the standard deviation of the population and the standard deviation √ of the sample means, and verify that means ≈ / n for the appropriate choice of n. 3. A recruiter for a government agency wants his candidates to be successful on a test administered by the agency to prospective employees. He would like his group of thirty recruits to score at least 700 out of the possible 950. Over the time this test has been administered, the average score has been 670, with a standard deviation of 105.

27.6

1185

(a) Estimate the probability that all thirty of the recruits will score in the 700–750 range. (b) Estimate the probability that all the recruits will score at least 670. (c) The recruiter gets a bonus if seven of the candidates score 800 or higher. Should the recruiter look forward to a bonus? 4. A uniform maker has a large order from a new branch of the military. On average, the pants legs are 31 inches, with a standard deviation of 32 inches. Suppose six new recruits are chosen at random. (a) Find the probability that their mean length is at least 29 inches. (b) Find the probability that their length is in the 29–33 inch range. (c) Find the probability that their length is no greater than 32 inches. 5. A state uses standard braces for its highway bridges. The braces can handle a mean pressure of 5000 pounds per square inch, with a standard deviation of 800. Suppose twenty braces are tested at random throughout the state. (a) What is the probability that their mean pressure tolerance is at least 4800 pounds per square inch? (b) What is the probability that their mean pressure tolerance falls in the 4700–5100 range? (c) What is the probability that their mean pressure tolerance is not less than 4500 pounds per square inch? 6. A product coming off an assembly line has been found to weigh on average 17 pounds, with a standard deviation of 18. Thirty samples are chosen at random from the line. (a) What is the probability that these samples will weigh on average between 168 and 172 pounds? (b) What is the probability that they weigh on average at least 167 pounds? (c) Find the probability that they do not weigh on average less than 172 pounds.

Confidence Intervals and Estimating Population Proportion If we somehow know that exactly five of every thousand cars coming out of a plant are defective, then there is no need to take samples, and we do not need statistics. The need for statistical analysis arises when we do not have perfect knowledge of the outcome of some

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process, as is usually the case in real-world scenarios. In this section we will focus on how to estimate (make an educated guess about) certain statistics, such as the mean and standard deviation of a population. We also want to be able to determine a confidence interval, which is a set of values measuring, in a sense we will define, ranges of accuracy to be expected of our estimates. Throughout this section we will have some population proportion p. We are examining some question or process which has exactly two outcomes, called success and failure. We form random samples of size n, and consider the sample proportions , p = x/n, where n is the size of the sample, and x is the number of successes in the sample. This number is between 0 and 1 and may be thought of as a probability. The complementary probability is , q = 1 −, p. As a typical example, we might draw samples of 900 trout from a stream over a day, and find that 47% are male. The sample size here is n = 900, and the sample proportion is , p = 047. We also assume that there is a fixed number of independent trials, and that the resulting distribution can be approximated by a normal distribution. In making estimates of various statistics or numbers associated with the population, we will use the notion of a confidence interval, which is an interval of values used to estimate the statistic. The corresponding confidence level of the estimate is the ratio of the times that this interval actually does contain the statistic of interest, taken over a large number of repetitions or samples. This confidence level is routinely denoted as 1 − . The most common confidence levels, expressed as percents, are 90%, 95%, or 99%. Below 90% is often unacceptable as a worthwhile measure of accuracy of the estimate. Associate with 90% the probability 09, so  = 01 and 1 −  = 09. Similarly, for 95%, use  = 005, so 1 −  = 095. And for 99%, use  = 001, so 1 −  = 099. If, for example, we are using 95% as a confidence level to estimate the population proportion p, and this determines (in a way we will describe) an interval a < p < b, then we say we are 95% confident that the true value of p, as estimated, falls between a and b. When these percentages are written decimals, we may think of them as probabilities, or areas bounded by the standard bell curve, which is the graph of 1 2 y = √ e−z /2  2 Part of this graph is shown in Figure 27.27. For a given , there is a point Z such that the area under the graph to the right of Z and to the left of −Z is , hence the area between −Z and Z is 1 − . This number Z, which depends on , is denoted z/2 . The use of /2 in the subscript is a reminder that the bell graph has two tails (to the left of −z/2 and to the right of z/2 ), having total area , so the area under just one tail is /2. The number z/2 is called a critical value, separating those sample proportions (between −z/2 and z/2 ) that are likely from those that are unlikely (in the right or left tails of the bell curve). Table 27.8 gives the critical values corresponding to confidence levels of 85%, 90%, 95%, and 99%. These percents are written as decimals 085, 090, 095, and 099 to have interpretations as areas. For example, the area under the standard bell curve between −z/2 and z/2 for  = 015 is 1
144 −z2 /2 e dz = 085013 √ 2 −144 corresponding to 85% of the sample proportions falling in this interval. We may also think of 085 as the probability that a sample proportion falls between −144 and 144 for this standard curve.

27.6 Confidence Intervals and Estimating Population Proportion

1187

0.4

0.3

0.2

0.1 Area α

Area α

2

2

0 –zα/2

zα/2

FIGURE 27.27 Critical point z/2 for the standard bell curve.

TA B L E 2 7 . 8 Confidence Level (As a Decimal) 085 090 095 099

 015 010 005 001

z/2 144 1645 196 2575

From previous results, we can write  = np and  =

√

npq

Since these results are for n trials, we can write the mean of the sample proportions as np =p = n and the standard deviations of the sample proportions as  √ npq pq =  = n n In general, we do not know p (or q), and we want to estimate p and , q.  these values from , Replace p by , p and q by , q in the expression for  to obtain , p, q /n. The probability is 1 − 

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that a sample proportion will differ  from the actual population proportion (which we want p, q /n. This leads us to define the maximum error of the to estimate) by no more than z/2 , estimate to be  , p, q  (27.6)  = z/2 n or  , p
1 −, p  = z/2  n The confidence level associated with this  is , p− < p 0279, there is a 095 probability that there is a significant linear correlation. If c is less than or equal to the critical value, then there is only a 005 probability of a linear correlation between the xi ’s and yi ’s. Notice that, as the number of sample points increases, we can infer significant linear correlations with smaller values of c, while, if n is small (and therefore provide less information), it takes a larger c to draw this inference.

TA B L E 2 7 . 9 n 10 15 20 25 30 35 40 45 50 60 80 100

Critical value of c for  = 005 0632 0514 0444 0396 0361 0335 0312 0294 0279 0254 0220 0196

27.8 Correlation and Regression

1197

EXAMPLE 27.28

The following table gives gold medal heights in the pole vault for the Olympics, from some years extending from 1932 through 1996. Heights are in feet. Year 1932 1936 1948 1952 1956 1960 1964 1968

Height 1415 1427 1410 1492 1496 1542 1673 1771

Year 1972 1976 1980 1984 1988 1992 1996

Height 1804 1804 1896 1885 1977 1902 1942

We will investigate this data set for a significant linear correlation, or absence of one. A scatter graph of the data is shown in Figure 27.33. The x-values are the years and the y-values are the heights. We find that x = 19669333, y = 169573, sx = 196813, and sy = 21138. A routine computation yields c = 09546. From Table 27.9 with n = 15, because c > 0514, we infer that there is a 95 probability that there is a significant linear correlation between x- and y-values. This is intuitively apparent from Figure 27.33.

25 20 15 10 5 0 1920

1930

1940

1950

1960

1970

1980

1990

2000

FIGURE 27.33 Olympic gold medal vaulting heights.

EXAMPLE 27.29

Consider the following data set: xi yi

0 1

05 212

10 624

15 244

20 1158

25 372

30 104

35 2563

40 263

Figure 27.34 shows a scatter graph of this data. There does not appear to be a significant linear correlation. To check this intuition, we find that c = 03884153476. Table 27.9 does not have

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30 25 20 15 10 5 0 0 FIGURE 27.34

1

2

3

4

5

Data with no significant linear correlation.

a critical value corresponding to n = 9, so we will use the value for n = 10 as a best estimate. This critical value is 0632. Since c < 0632, there is only a 005 probability that there is a significant linear correlation between the x-values and the y-values. Correlation must not be confused with causality. Just because two sets of data have a significant linear correlation, does not imply that the x-values cause the y-values. There may or may not be such a cause and effect relationship. In the case of pole vaulting heights, it is hard to believe that the year is a cause of the height achieved in the pole vault. The fact that the heights increase with passing years is due to a complex set of factors, probably including conditioning of the athletes, improvements in vaulting technique, and changes in materials from which vaulting poles were made. To understand the idea of regression, consider again the data shown in Figure 27.33. We found that there was a significant linear correlation between x-values and y-values. This suggests that the data, plotted as points in the plane, form an “approximate” straight line. It is possible to determine a straight line L which is a best fit to the data. This line is called the regression line for the data. It is defined to be the best least squares fit to the data, in the sense of minimizing the sum of the squares of the vertical distances between the yi s and the respective y-values on the line corresponding to the xi s. If we write the equation of the regression line as y = a + bx, then the standard least squares method yields     
ni=1 yi  ni=1 xi2 −
ni=1 xi 
ni=1 xi yi  a= (27.10)    n ni=1 xi2 −
ni=1 xi 2 and b=

n

n

n n i=1 xi yi −
i=1 xi 
i=1 yi      n ni=1 xi2 −
ni=1 xi 2

(27.11)

We can also compute a from the expression a = y − bx if we have computed the means of the xi ’s and the yi ’s.

(27.12)

27.8 Correlation and Regression

1199

EXAMPLE 27.30

We will determine the regression line for the heights of Olympic pole vaults. Use equations (27.11) and (27.12) to compute a and b, obtaining a = −185637 and b = 0103 The regression line for this data has equation y = −185637 + 0103x Figure 27.35 again shows the data (as in Figure 27.33), along with the regression line. The data ends at 1996. If we use the regression line to try to project ahead to 2008, we get y = −185637 + 0103
2008 = 21187 ft This would require an extraordinary effort. But we learn a more dramatic lesson if we use the line to project to the 3000 Olympics. Now we get y = −185637 + 0103
3000 = 12336 ft No one would believe this. The point is that the regression line in this example has positive slope, and it will return as large as value of y as we want by putting in a sufficiently large value of x. However, there are clearly limits to how high a human being can vault. We do not know what this limit is, but 124 feet seems too high. As this example shows, we must not use a regression line to attempt to project conclusions about y-values that are “too far” outside the data set from which the line was derived. The further out of touch we get from the data, the less reliable our projected conclusions are. There is a very useful interpretation that can be made for the square of the linear correlation coefficient c. Suppose we have sample points
x1  y1     
xn  yn , and there is a significant linear correlation between the xi ’s and the yi ’s. The regression line is y = a + bx, with a and b from equations (27.10) or (27.12), and (27.11). For a given xi , write
yr i = a+bxi to distinguish between yi and the y-coordinate of the point on the regression line corresponding to xi . 25

20

15

10

5

0 1920

1930

1940

1950

1960

1970

1980

FIGURE 27.35 Olympic vaulting heights and the regression line.

1990

2000

1200

CHAPTER 27

Statistics

If we think of the regression line as defining a sense of how the data “should” arrange itself, we can look to the regression line to “explain” the y-values. That is, think of
yr i as the number that should go with xi , and yi as the number that actually does go with xi in the sample. This makes it reasonable to attempt to compute the ratio of differences of the data from the mean that are explained by the regression line, to the total differences of the data from the mean. We can quantify these ideas as follows. Interpret n 


yi − y

i=1

as the total deviation of the yi s from their mean. Further, n 



yr i − y

i=1

is the deviation of regression line y-values from the mean of the yi ’s. Think of this as a deviation explained by the regression line. The ratio we are interested in, which we will temporarily call k, is therefore n i=1

yr i − y  k=  n i=1
yi − y Next, n 


yi −
yr i 

i=1

is the deviation of the regression line values from the actual y-values. Think of this as the deviation for which the regression line presents no rationale. Now observe that n n n   
yi − y =

yr i − y +
yi −
yr i   i=1

Therefore,

i=1

i=1

n

i=1

yr i − y n i=1

yr i − y + i=1
yi −
yr i 

k = n =

differences of yi s from the mean explained by the regression line  total differences of yi ’s from the mean

Now comes what may be a surprise. If the algebra is carried out, we find that k = c2 . Thus c enables us to detect a significant linear correlation, and c2 gives a measure of how well the deviations of the yi ’s from the norm are explained by the regression line.

EXAMPLE 27.31

A state compiles statistics on deer hunters and deaths in the deer population each year. For a period of ten years, these are given by the following table: Number of Hunters Deaths in Deer Population

168 28

172 30

194 36

204 39

230 42

270 50

295 46

320 57

390 61

402 62

We find that x = 25550, y = 4510, sx = 97961, and sy = 1228. A routine computation yields c = 08488. With n = 10 in Table 27.9, the critical value is 0632. Because c > 0632, there is

27.8 Correlation and Regression

1201

70

Number of deer deaths

60 50 40 30 20 10 0 0

100

200

300

400

500

Number of hunters FIGURE 27.36 Hunters vs. deer deaths.

a 95 probability that there is a significant linear correlation between x- and y-values. A scatter plot of these values is given in Figure 27.36. Now we ask: How many of the deaths in the deer population are explained by the number of hunters? We expect some deer to be killed by hunters, and some to die from other causes. Compute c2 =
084882 = 07201 We conclude that 72% of the deer deaths can be explained by the regression line, which means that 72% can be explained by the presence of the hunters. The other 28% must be accounted for by other factors. Finally, we would like to develop a kind of confidence interval for predictions based on a regression line. In the present context, such an interval will be called a prediction interval. Begin with the spread of the error of the estimate, which is analogous to standard deviation in the sense that it measures the spread of the sample points about the regression line. This spread of the error of the estimate is denoted Sr and is defined by  n 2 i=1
yi −
yr i   Sr = n−2 It is routine to verify that 0 Sr =

n

2 i=1 yi − a

n

i=1 yi − b

n−2

n

i=1 xi yi



with a and b given by equations (27.10) and (27.11) (or (27.12)), respectively.

1202

CHAPTER 27

Statistics

Next, analogous to the process developed for confidence intervals, define the maximum error r for a given x-value (which is in the x-population, but need not be one of the actual sample values xi ) to be 0 n
x − x2 1 r = t/2 Sr 1 + + n 2  n  n n i=1 xi −
i=1 xi 2 Here t/2 is a critical point of the Student t-distribution corresponding to a confidence level . We have usually taken  = 005 to have a 95% confidence level in the conclusion. The number t/2 must be determined from a table. The difference between the current and previous use is that previously, in estimating population mean, we had n − 1 degrees of freedom, and now we have n − 2. This information must be used in order to reference the value of t/2 from a table of critical values for the t-distribution. Now, for a given x, we define the prediction interval for the corresponding y to be yx − r < y < yx + r  where yx = a + bx, the y-value given by the equation of the regression line for the given value of x.

EXAMPLE 27.32

Consider again the hunters and deaths in the deer population given over a ten year period. Suppose we want to estimate how many deer kills would correspond to x = 350 hunters. It is routine to compute the equation of the regression line, obtaining y = 84509 + 01386x. We have stated that x = 25550. We need to compute
yr i = a + bxi for i = 1     10. These values are: xi 168 172 194 204 230 270 295 320 390 402
yr i = a + bxi 31.7357 32.2901 35.3393 36.7253 40.3289 45.8729 49.3379 52.8029 62.5049 64.1681

Compute 10 
yi −
yr i 2 = 803589 i=1

and

0 Sr =

Then

10

i=1
yi −
yr i 

8 0

r =
2306
31694 1 +

2

= 31694

10
350 − 25550 1 + 10 10
765989 −
26452

= 76703 Finally, compute yx = a + bx = 84509 + 01386
350 = 59961, the y-value on the regression line corresponding to x = 350. The prediction interval is yx − r < y < yx + r , or 59961 − 76703 < y < 59961 + 76703

27.8 Correlation and Regression

1203

In the present context, we will use two decimal places and write the prediction interval as 5229 < y < 6763 The meaning of this interval is that, for a hunter population of 350, the predicted (by the regression line) number of deer deaths is 5996 (or 60) , and we have a confidence level of 95% (because of the way we chose the critical value t/2 ) that the actual number is between 52 and 68. Again, we can change the confidence level from 95% to another value by making the appropriate change in the choice of t/2 .

SECTION 27.8

PROBLEMS

For each of Problems 1 through 8, draw a scatter plot of the data, compute c and determine whether or not there is a significant linear correlation between the xi ’s and the yi ’s. If there is not, the problem is completed. If there is, determine the equation of the regression line, and, for each given x-value, determine the y-value predicted by the regression line, as well as the corresponding prediction interval with a 95% confidence level.

7. 0 −31

2.5 1.22

6.1 7.2

8.3 10.37

9.8 12.9

11.6 15.85

15.3 20.7

19.5 28.19

22.9 32.7

25.5 38.6

29.6 43.92

32.2 46.8

36.2 56.2

40.8 61.9

17.2 25.2

x = −4 111 301, and 667. 8.

2. 0.2 0.9 1.4 1.9 2.6 3.4 4.2 6.8 10 15.9 2.8 0.93 −061 −21 −39 −65 −89 −161 −272 −4241

xi yi

0 3

1.3 0.5

1.9 −084

5.4 7.4 12.9 −74 −1232 −248

2.3 −174 15.3 −276

2.5 −190 19.4 −359

2.8 −273 22.5 −437

3.1 4.8 −30 −69 25.9 −485

x = 15 43 13, and 32.

x = 52 86, and 178. 3. xi yi

x = 35 47, and 73

0.98 1.3 2.1 2.8 3.6 5.2 5.8 7.3 9.7 13.6 3.91 4.10 4.81 5.07 6.21 7.15 7.66 9.04 10.84 14.26

x = 04 62, and 151.

xi yi

xi 3.7 3.9 4.2 4.6 4.8 5.2 5.7 6.3 6.5 6.9 yi −12 3.9 −56 3.8 −129 4.2 −272 4.8 −186 15.9

xi yi

1. xi yi

6.

1.9 2.5 5.4 8.91 10.5 17.3 18.9 21.7 25.3 −14 4.2 3.15 −56 10.7 11.2 −68 −153 2.7

34.7 9.4

x = 0 92, and 388.

9. A basketball coach is attempting to evaluate the impact of an expensive player on the team’s won/loss record. He considers the following table, which gives number of games per season in which the player has seen at least twenty five minutes of floor time, together with wins over the season.

4. xi 4.8 5.8 7.3 9.5 15.3 19.5 27.4 32.6 36.3 41.9 yi 34.8 41.1 47.2 57.1 76.9 98.6 133.8 148.4 164.9 197.5

xi : Games yi : Wins

62 47

58 46

59 47

73 56

69 55

75 62

72 55

81 63

46 38

79 64

(a) Show that there is a significant linear correlation between games in which the player participated, and team wins.

x = 29 29 and 46. 5.

(b) Write the equation of the regression line. xi yi

3 1

34 13

38 19

42 228

59 382

75 542

x = 06 30 390, and 421.

123 90

218 191

298 249

353 287

(c) Estimate the percentage of wins (in games the player played in) that are explained by this player’s floor presence.

CHAPTER 27

1204

Statistics

(d) Predict with a 95% confidence level how many games the team will win next season if this player is in 85 games. Write the prediction interval for this estimate.

(a) Show that there is a significant linear correlation between number of visits and the number of thefts.

10. A drug is being tested. One purpose of the test is to demonstrate the efficacy of the drug, but another is to measure incidence of a certain side effect (lowered body temperature). The following table gives numbers of randomly chosen patients in ten states, together with the number exhibiting the side effect:

(c) Estimate the percentage of thefts that are explained by this person’s number of visits to the city each year.

xi : Patients yi : Side Effect

5,000 160

5,500 167

4,200 135

10,000 305

6,900 221

8,800 265

12,500 394

6,000 190

8,000 235

5,700 183

(a) Show that there is a significant linear correlation between number of patients and incidents of this side effect.

(b) Determine the equation of the regression line.

(d) Predict with a 95% confidence level how many thefts are likely to occur if this person visits New York 25 times next year. Write the prediction interval for this estimate. 12. A city is experimenting with an additive to the water with the objective of decreasing tooth decay. Over a ten year period, records have been kept on pounds of the chemical put into the water system and the number of children twelve or under who were cavity free that year.

(b) Determine the equation of the regression line. (c) Estimate the percentage of patients that will have this side effect in a patient population of 17 000. (d) Predict with a 95% confidence level how many patients out of a population of 17 500 will exhibit this side effect. Write the prediction interval for this estimate. 11. Over a period of ten years, the number of times a certain individual has visited New York each year has been recorded. Records have also been kept of the number of incidents of major jewelry thefts in the city for each year in this period. The numbers are as follows: xi : Visits yi : Thefts

17 6

12 5

15 6

22 8

14 5

9 3

18 7

21 8

10 5

13 5

xi : Additive (pounds) yi : Cavity free

410 3550

500 4412

475 4119

390 3420

470 4090

525 4598

480 4203

510 4452

590 5143

600 5259

(a) Show that there is a significant linear correlation between number of pounds of the chemical, and the number of children who were cavity free. (b) Determine the equation of the regression line. (c) Estimate the percentage of children who were cavity free that is explained by the chemical added to the water supply. (d) Predict with a 95% confidence level how many children are likely to be cavity free if 650 pounds are used. Write the prediction interval for this estimate.

Answers and Solutions to Selected Problems CHAPTER 1 Section 1.1 √ 1. Yes, since 2
= 2
x − 1

 √

1

2 x−1

 = 1 for x > 1

3. Yes

5. Yes

d 2
y + xy − 2x2 − 3x − 2y =
y − 4x − 3 +
2y + x − 2y
= 0 dx 13. y = 3 − e−x 15. y = 2 sin2
x − 2 7.

17. Direction field for y
= x + y; solution satisfying y
2 = 2

19. Direction field for y
= xy; solution satisfying y
0 = 2

y


4

y

4

4

2

2


2

2

4

x


4


2

2


2


2


4


4

21. Direction field for y
= sin
y; solution satisfying y
1 = /2

y

4

4

2

2


2

x

23. Direction field for y
= y sin
x − 3x2 ; solution satisfying y
0 = 1

y


4

4

2

4

x


4


2

2


2


2


4


4

4

x

A1

Answers and Solutions to Selected Problems

A2

25. Direction field for y
− y cos
x = 1 − x2 ; solution satisfying y
2 = 2

y 4 2
4


2

2

4

x


2
4 27. Hint: The lineal element to the solution curve through
x0  z has equation y = q
x0  − p
x0 z
x − x0  + z

Section 1.2 1 ; also y = 0 7. sec
y = Ax 9. Not separable 1 − cx 1 2 2 11. 2 y − y + ln
y + 1 = ln
x − 2 13.
ln
y = 3x2 − 3 15. 3y sin
3y + cos
3y = 9x2 − 5  3   √ 4 1/3 1/3  17. 45 F 19. V = V0 + k t ; V0 = initial volume 21. 8.57 kg 23. 21 e−6 3 √ 25. 3888 2 s, or approximately 91 min, 39 s 27. (a) 576 s (b) 489 s (approximate) 29. No 1. 2x2 = y3 + c

3. Not separable

5. y =

Section 1.3 1. y = cx3 + 2x3 ln x 3. y = 21 x − 41 + ce−2x 5. y = 4x2 + 4x + 2 + ce2x
x−1 ex dx + ce−x 9. y = x2 − x − 2 11. y = x + 1 + 4
x + 1−2 13. y = 23 e4x − 11 7. y = e−x ex 3 x2 15. y = −2x2 + cx 17. A1
t = 50 − 30e−t/20  A2
t = 75 + 90e−t/20 − 75e−t/30  A2
t has its minimum value of 5450 81   pounds at 60 ln 95 minutes.

Section 1.4 1. 2xy2 + exy + y2 = c

5. y3 + xy + ln x = c 7. cosh
x sinh
y = c  9. 3xy4 − x = 47 11. x sin
2y − x = 13. Not exact 15.  = −3 x2 y3 − 3xy − 3y2 = c 24



 17.
 + c = = M and
 + c = =N

x

x

y

y 3. Not exact

Section 1.5 1. 3.

5. 11. 19.



N M − is independent of x.

x

y

N 1

M = 1 = −1, so this equation is not exact. (b) 
x = 2 (c) 
y = (a)

y

x x (d)
x y = xa yb for all a and b satisfying a + b = −2 1  x2 y = c or y = −1 e3y  xe3y − ey = c 7. x2 y x4 y2 + 2x3 y3 = c 9. y+1 1 e−3x y−4  y3 − 1 = ky3 e3x 13.  y = 4 − ln x 15. x x2
y3 − 2 = −9 17. x




cM = c
M = c
N =
cN ex  ex sin
x − y = 21 21.

y

y

x

x 1 M



1 y2

1 2  y = 4e−x /3 y

Answers and Solutions to Selected Problems

A3

Section 1.6 x 2 3. y = 1/
1 + ex /2  5. y ln y − x = cy 7. xy − x2 − y2 = c c − ln x  4/7 2ex ce − br ar − dc 2 13. y = 2x 15. h = k= 11. y = 2 + 2 9. y = x−1 c − 75 x−5/4 cx − 1 ce − 1 ae − bd ae − bd 17. 3
x − 22 − 2
x − 2
y + 3 −
y + 32 = c 19.
2x + y − 32 = c
y − x + 3 21.
x − y + 32 = 2x + c 23. 3
x − 2y − 8 ln x − 2y + 4 = x + c 25. Let the dog start at
A 0, and suppose the man walks upward into the upper half-plane. The dog’s path is along the √ √ 1 2 graph of y = − A A − x + √
A − x3/2 + A, and the dog catches the man at
A 2A/3 at time t = 2A/3v 3 3 A units. v a 27. (a) r = a − , a spiral (b) revolution  2v    √ a + v2 + a2 2 1 a  2 2 2 v + a  + ln (c) Distance = 2 v2 v 1. y = x +

Section 1.7

 √ 8 1358 , about 18.91 ft/s 3. Velocity = 8 30, about 43.82 ft/s 9 3 √ √ √ Velocity = 12 5, about 26.84 ft/s; time = 23 ln
6 + 35, about 2.15 s √
√ 3 40 x Time = dx, about 1.7117 seconds 9. Velocity = 2 210, about 28.98 ft/s √ 8 10 x3 − 100 Max height = 342.25 ft; object hits the ground at 148 ft/s, 4.75 s after drop v
t = 32 − 32e−t for 0 ≤ t ≤ 4 v
t = 8
1 + ke−8t /
1 − ke−8t  ft/s for t ≥ 4, with k = e32
3e4 − 4/
5e4 − 4 limt→ v
t = 8 ft/s; s
t = 32
t + e−t − 1 for 0 ≤ t ≤ 4 s
t = 8t + 2 ln
1 − ke−8t  + 64 + 32e−4 − 2 ln
2e4 /
5e4 − 4 for t ≥ 4 √ 17.5 ft/s 17. t = 2 R/g, where R = radius of the circle VC = 76 V when t = ln
20/2; i
21 ln
20 = 16 A 21. i1
0+ = i2
0+ = 3/20 amp i3
0+ = 0   001 EC −t/RC (a) q
t = EC +
q0 − ECe (b) EC (d) −RC ln q0 − EC 3 y = c − 4 ln x 27. x2 + 2y2 − 4y = c 29. y2
ln
y2  − 1 + 2x2 = c

1. Velocity = 5. 7. 11. 13.

15. 19. 23. 25.

Section 1.8 1 1 1
−1n+1 n x (c) yn = 1 + x − x2 + x3 − x4 + · · · + 2 6 24 n! n n+1 
−1 n 1 1 1
−1 x = 1 + x − x2 + x3 − x4 + · · · + xn + · · · (d) y = 2 − n! 2 6 24 n! n=0

7. (b) y = 2 − e−x

9. (b) y = 73 + 23 x3 (c) y0 = 3 y1 = 73 + 23 x3 = y2 = y3 = · · · = yn (d) y = 3 + 2
x − 1 + 2
x − 12 + 23
x − 13 = y1

CHAPTER 2 Section 2.2 1. y = cosh
2x

3. y =

12 −3x e 5

− 75 e−8x

5. y = 2x4 − 4x4 ln x

4 6 7. y1 and y2 are solutions of x2 y

− 4xy
+ 6y = 0. We must write this equation as y

− y
+ 2 y = 0 to have the form x x to which the theorem applies, and then we must have x = 0. On an interval not containing 0, the Wronskian of these solutions is nonzero.

A4

Answers and Solutions to Selected Problems

9. y

− y
− 2y = 0 has solutions y1 = e−x and y2 = e2x , but y1 y2 = ex is not a solution. Many other examples also work. 11. Hint: Recall that at a relative extremum x0  y

x0  = 0. 13. Hint: Consider W
x0 .

Section 2.3 1. y = c1 cos
2x + c2 sin
2x 3. y = c1 e5x + c2 xe5x 5. y = c1 x2 + c2 x2 ln x 7. y = c1 x4 + c2 x−2     cos
x sin
x 9. y = c1 + c2 11. y = c1 e−ax + c2 xe−ax √ √ x x 13. (a) y4 = c1 x + c2 (b)
y − 1ey = c1 x + c2 or y = c3 (c) y = c1 ec1 x /
c2 − ec1 x  or y = 1/
c3 − x (d) y = ln  sec
x + c1  + c2 (e) y = ln c1 x + c2 

Section 2.4 1. y = c1 e−2x + c2 e3x 7. 13. 17. 21.

3. y = e−3x c1 + c2 x 5. y = e−5x c1 cos
x + c2 sin
x √ √ √ √

c1 cos
3 7x/2 + c2 sin
3 7x/2 9. y = e7x
c1 + c2 x 11. y = e−2x
c1 cos
5x + c2 sin
5x y=e y = 5 − 2e−3x 15. y
x = 0 y = 97 e3
x−2 + 57 e−4
x−2 19. y = ex−1
29 − 17x  √ √  y = e
x+2/2 cos 215
x + 2 + √515 sin 215
x + 2 −3x/2

1 ax e
d − ac + cex +
ac − d + ce−x  2 (c) lim→0 
x = 
x (use l’Hôpital’s rule)

23. (a) 
x = eax c +
d − acx

(b) 
x =

Section 2.5 1. y = c1 x2 + c2 x−3 7. 13. 17. 19.

5. y = c1 x4 + c2 x−4 √ √ y = c1 x + c2 x 9. y = x
c1 + c2 ln
x 11. y = x c1 cos
39 ln
x/2 + c2 sin
39 ln
x/2 y = x−2
3 cos
4 ln
−x − 2 sin
4 ln
−x 15. y = −3 + 2x2 y = −x−3 cos
2 ln
−x y = −4x−12
1 + 12 ln
x 21. y = 11 x2 + 17 x−2 4 4 −2

3. y = c1 cos
2 ln
x + c2 sin
2 ln
x

−3

−12

3/2

Section 2.6 1. y = c1 cos
x + c2 sin
x − cos
x ln  sec
x + tan
x 3. y = c1 cos
3x + c2 sin
3x + 4x sin
3x + 43 ln  cos
3x cos
3x −x

5. y = c1 ex + c2 e2x − e2x cos
e−x 

7. y = c1 e + c2 e − x + x − 4 9. y = e c1 cos
3x + c2 sin
3x + 2x2 + x − 1 11. y = c1 e2x + c2 e4x + ex 13. y = c1 ex + c2 e2x + 3 cos
x + sin
x 15. y = e2x c1 cos
3x + c2 sin
3x + 13 e2x − 21 e3x 2x

2

x

17. y = c1 e2x + c2 e−x + 13 xe2x 19. y = c1 e2x + c2 e−3x − 16 x − 361

21. y = c1 x2 + c2 x4 + x

23. y = c1 cos
2 ln
x + c2 sin
2 ln
x − 41 cos
2 ln
x ln
x

25. y = c1 e2x + c2 xe2x + e3x − 41

27. y = 47 e2x − 43 e−2x − 47 xe2x − 41 x 19 −6x 29. y = 38 e−2x − 120 e + 15 e−x + 127

31. y = 2e4x + 2e−2x − 2e−x − e2x

33. y = − 17 e2x + 55 e3x + 521 cos
2x − 525 sin
2x 4 13  √ √ 1 4
x+1 35. y = 11
55 − e cosh
14
x + 1 + √114
5e − 198 sinh
14
x + 1 + 111 e−x e 37. y = 4e−x − sin2
x − 2

39. y = 2x3 + x−2 − 2x2

41. y = x − x2 + 3 cos
ln
x + sin
ln
x

Answers and Solutions to Selected Problems

A5

Section 2.7 √ √ 1. y = e−2t 5 cosh
2t + √102 sinh
2t √ y = √52 e−2t sinh
2t
graphed below

3. y = 25 e−t 2 cos
2t + sin
2t y = 25 e−t sin
2t
graphed below

y

y 5

5

4

4 3

3

y(0) = 5

2

2 y(0) = 0

1

1 t 0

2

4

6

8

10

2

0

4


1

12

√ A 5. y = √ e−2t sinh
2t
graphed below 2

y A  10

1.5

1.5 1.0

6

0.5

3

0
4


0.5
1.0

9. y =

1.0

A1

2

4

0.5 6

8

t

10


1.0

A 
7

A −t e sin
2t
graphed below 2

2 1 2


1

0
0.5

y

0

t

8

7. y = Ate−2t
graphed below

y 2.0

6

4

6

8

t

2

4

6

8

t

Answers and Solutions to Selected Problems

A6

11. y = 601
57e−8t − 52e−3t  down positive 13. At most once; no condition on only y
0 is enough, as one needs to specify y

0 also. 15. Increasing C decreases the frequency 17. Hint: Show that
m/k2 =
d/a2    √ √ 1 √ sinh
7t − 28 cos
3t + 72 sin
3t 21.
a y = 373 e−3t 2266 cosh
7t + 6582 7    √ √ 1 √ sinh
7t − 28 cos
3t + 72 sin
3t
b y = 373 e−3t 28 cosh
7t + 2106 7 These functions are graphed below.

19.

A t sin
0 t 2m0

y 6 5 4 3 2 1 t

0

2

4

6

8

10

12

14

  √ √ e−t/2 98 cos
11t/2 + √7411 sin
11t/2 − 8 cos
3t + 4 sin
3t    √ √ 1 e−t/2 8 cos
11t/2 + √164 sin
11t/2 − 8 cos
3t + 4 sin
3t
b y = 15 11 These functions are graphed below.

23.
a y =

1 15



y 6 4 2 2

4

6

8

10

12

14 t

0
2

25. i
t = 015e−00625t − 54 × 10−7 e−333327t + 0015 cos
20t − 000043 sin
20t 27. i
t = 00001633e−t + 000161e−03177t + 0000023e−t cos
6t − 0000183e−t sin
6t

CHAPTER 3 Section 3.1 1.

1 1 − s−1 s+1

3.

13. 2 cos
4t − 45 sin
4t

16s
s2 + 42

5.

15. 3e7t + t

1 s − s2 s2 + 25

7.

17. e4t − 6te4t

2 6 8 16 3 s 9. 4 − 2 + 2 + + s3 s2 s s s s + 16 5 23. s
1 + e−3s 

11. −2e−16t

Answers and Solutions to Selected Problems

A7

25. From the graph, f
t = 0 if 0 < t ≤ 5, and f
t = 5 if 5 < t ≤ 10, and f
t = 0 if 10 < t ≤ 25. Further, 5e−5s
1 − e−5s  . f
t + 25 = f
t, so f is periodic of period T = 25. Thus £ f
s = s
1 − e−25s  E 1 29. f
t = h if 0 < t ≤ a and f
t = 0 if a < t ≤ 2a. Further, f
t + 2a = f
t, so f is periodic 27. 2 s + 2 1 − e−s/ h . of period 2a, and L f
s = s
1 + e−as 

Section 3.2 1. y = 41 − 13 e−4t 4

4 −4t 4 3. y = − 17 e + 17 cos
t + 171 sin
t

7. y =

+

22 2t e 25

−

13 2t te 5

3 25

cos
t −

4 25

sin
t

9. y =

1 16

+

5. y = − 41 + 21 t + 17 e2t 4 1 t − 33 16 16

15 cos
4t + 64 sin
4t

Section 3.3 1 6 3 2 s 1 3.
1 − e−7s  + 2 cos
7e−7s − 2 sin
7e−7s − +
s + 24
s + 22 s + 2 s s +1 s +1   2 s 2 s 1 1 11 4 1 1 − 9. + − − e−2s 5. 2 − e−3s − 2 e−3s 7. + s s s s + 1
s + 13
s + 12 + 1 s2 + 1 s s2 + 1 s2 + 1   s2 + 4s − 5 1 2 24 1 15 −16s 4 1 e 11. 2 13. 2 − − 2 + 15. + + 17. e2t sin
t
s + 4s + 132 s s s s
s + 55
s + 53
s + 52 √ √ √ 19. cos
3
t − 2H
t − 2 21. √12 e−3t sinh
2t 23. e−3t cosh
2 2t − 2√1 2 e−3t sinh
2 2t 1.

25. 29. 31.

1

1 − cos
4
t − 21H
t − 21 27. y = cos
2t + 43 1 − cos
2
t − 4H
t − 4 16 √ y = − 41 + 121 e2
t−6 + 16 e−
t−6 cos
3
t − 4H
t − 6 y = − 41 + 25 et − 203 cos
2t − 15 sin
2t + − 41 + 25 et−5 + 203 cos
2
t − 5 − 15 sin
2
t − 5H
t − 5

33. Eout = 5e−4t + 10
1 − e−4
t−5 H
t − 5

35. i
t =

k k
1 − e−Rt/L  −
1 − e−R
t−5/L H
t − 5 R R

K K 37. £ KH
t − a − KH
t − b
s = e−as − e−bs s s  c−t  c−t    t−a t−a H
t − a + h − H
t − b − h H
t − c 39. £ h b−a c−b b−a c−b h
c − a e−bs h e−cs h e−as − + = b − a s2
c − b
b − a s2 c − b s2

Section 3.4 cos
at − cos
bt t sin
at if b2 = a2  if b2 = a2
b − a
b + a 2a   1 1 5. 4 1 − cos
at − 3 t sin
at 7. 21 − 21 e−2
t−4 H
t − 4 9. y
t = e3t ∗ f
t − e2t ∗ f
t a 2a 11. y
t = 41 e6t ∗ f
t − 41 e2t ∗ f
t + 2e6t − 5e2t 13. y
t = 13 sin
3t ∗ f
t − cos
3t + 13 sin
3t 1.

1

sinh
2t − sin
2t 16

3.

15. y
t = 43 et − 41 e2t − 121 e−2t − 13 et ∗ f
t + 41 e2t ∗ f
t + 121 e−2t ∗ f
t 17. f
t = 21 e−2t − 23 √ √ 19. f
t = cosh
t 21. f
t = 3 + 25 15et/2 sin
15t/2 23. f
t = 41 e−2t + 43 e−6t

Section 3.5 1. y = 3 e−2
t−2 − e−3
t−2 H
t − 2 − 4 e−2
t−5 − e−3
t−5 H
t − 5 3. y = 6
e−2t − e−t + te−t  5. 
t =
B + 9e−2t −
B + 6e−3t  
0 = 3 

0 = B 7. 3/ 9. 4 11. Eout = 10e−4
t−2 H
t − 2 − 10e−4
t−3 H
t − 3 13. 0 if t < a f
t − a if t ≥ a    m k 15. y
t = v0 sin t k m

A8

Answers and Solutions to Selected Problems

Section 3.6 1. x
t = −2 + 2et/2 − t y
t = −1 + et/2 − t

3. x
t = 49 + 13 t − 49 e3t/4  y
t = − 23 + 23 e3t/4

5. x
t = 43 − 43 e2t/3 + 21 t2 + 21 t y
t = − 23 e2t/3 + t + 23 7. x
t = e−t cos
t + t − 1 y
t = e−t sin
t + t2 − t 9. x
t = 1 − e−t − 2te−t  y
t = 1 − e−t 11. y1
t = 21 et + 21 e−t − 1 − t y2
t = − 41 t2 − 21 t y3
t = − 16 et + 16 e−t − 13 t 2 13. i1
t = 15
1 − 21 e−t/2  − 85

e−
t−4/2 − cos
2
t − 4 + 29 sin
2
t − 4H
t − 4 and 1 −t/2 2 −
t−4/2 + 85 e − cos
2
t − 4 − 4 sin
2
t − 4H
t − 4 i2
t = 10 e

 5 1 4 4 15. x1
t = 36 − 20 cos
2t − 45 cos
3t − 365 − 201 cos
2
t − 2 − 45 cos
3
t − 2 H
t − 2 and 

2 2 x2
t = 181 − 101 cos
2t + 45 cos
3t − 181 − 101 cos
2
t − 2 + 45 cos
3
t − 2 H
t − 2

17. m1 y1

= k
y2 − y1  m2 y2

= −k
y2 − y1  y1
0 = y1

0 = y2

0 = 0 y2
0 = d. Then
m1 s2 + kY1 − kY2 = 0 and 2
m2 s2 + kY2 − kY1 = m2 ds. Replace Y2 with m1 sk +k Y1 in the second equation to get kd Y1
s = . The quadratic term in the denominator indicates that the objects will m1 s s2 + k

m1 +  m2 /m1 m2  m1 m2 . oscillate with period 2 k
m1 + m2 

1 −
t−1 3 −
t−1/6 

 + 20 e 19. i1 = 10 e H
t − 1 i2 = − 101 e−
t−1 + 101 e−
t−1/6 H
t − 1 21. x1
t = 9e−t/100 + e−3t/50 + 3 e−
t−3/100 − e−3
t−3/50 H
t − 3 x2
t = 6e−t/100 − e−3t/50 + 2e−
t−3/100 + 3e−3
t−3/50 H
t − 3

Section 3.7 1. y = −1 + ce−2/t

3. y = 7t2

5. y = ct2 e−t

7. y = 4

9. y = 3t2 /2

CHAPTER 4 Section 4.1 1 6 1. y = −2 − 13 x3 + 121 x4 − 601 x5 − 120 x +···

3. y = 3 + 25
x − 12 + 56
x − 13 + 245
x − 14 + 16
x − 15 + · · ·

5. y = 7 + 3
x − 1 − 2
x − 12 −
x − 13 +
x − 14 + · · · 7. y = −3 + x + 4x2 + 76 x3 + 13 x4 + · · · 1 1 1 9. y = 1 + x − x2 + x3 + x4 + · · · 4 4 32 



31 6 379 8 5 6 43 8 1741 11. y = a0 1 − 21 x2 + 16 x4 − 720 x + 40320 x − 21 x2 − 18 x4 + 144 x − 5760 x + 1209600 x10 + · · ·

1 1 4 1 6 1 x − x + x8 + · · ·  13. y = x − 1 + A 1 − x2 + 2 4
2! 8
3! 16
4!  1 2n 2 x = x − 1 +
A + 1e−x /2 y = x − 1 + A
−1n n!2n n=0 an−1 1 6 1 1 for n ≥ 1
a1 = 0 y1 = 1 − 16 x3 + 180 15. an+2 = − x − 12960 x9 + 1710720 x12 − · · · 
a0 = 0
n + 1
n + 2 1 7 1 1 y2 = x − 121 x4 + 504 x − 45360 x10 + 7076160 x13 − · · ·    1 1 1 1 1 17. y = a0 1 + 4
5 x5 + 4
5
9
10 x10 + 4
5
9
10
14
15 x15 + · · · + a1 x + 5
6 x6 + 5
6
10
11 x11 +  1 1 1 1 x16 + · · · + 21 x2 + 2
6
7 x7 + 2
6
7
11
12 x12 + 2
6
7
11
12
16
17 x17 + · · · 5
6
10
11
15
16

1 1 1 1 6 1 7 1 1 7 5 x − x + · · · + a1 x − x2 + x3 − x4 + x −··· 19. y = a0 1 + x4 − x5 + 12 60 360 2520 2 6 24 120     1 1 1 1 1 1 3 6 9 4 7 21. y = a 1 − 2
3 x + 2
3
5
6 x − 2
3
5
6
8
9 x + · · · + b x − 3
4 x + 3
4
6
7 x −
3
4
6
7
9
10 x10 + · · · 3 3x3 1 4 23. y = 2 − 3e−x = −1 + 3x − x2 + − x +··· 2 3! 8 x4 + · · · 25. y = 15 + 15 e2x 32 sin
x − 6 cos
x = −1 + 4x + 11x2 + 343 x3 + 27 4

Answers and Solutions to Selected Problems Section 4.2 1 1. a0 arbitrary, a1 = 1 2a2 − a0 = −1 an = an−2 for n ≥ 3, n   1 1 1 1 1 1 x4 + 2
4
6 x6 + 2
4
6
8 x8 + · · · + x + 13 x3 + 3
5 x5 + 3
5
7 x7 + 3
5
7
9 x9 + · · · y = a0 +
a0 − 1 21 x2 + 2
4 1 3. a0 arbitrary, a1 + a0 = 0 2a2 + a1 = 1 an+1 =
−an + an−2  for n ≥ 2 n +  1 2 1 13 1 4 11 5 31 6

1 2 1 3 7 4 y = a0 1 − x + 2! x + 3! x − 4! x + · · · + 2! x − 3! x + 4! x + 5! x − 6! x + · · ·
n − 1 a for n ≥ 1 5. a0  a1 arbitrary, a2 = 21
3 − a0 , an+2 =
n + 1
n + 2 n   y = a0 + a1 x +
3 − a0  2!1 x2 + 4!1 x4 + 6!3 x6 + 3
5 x8 + 3
5
7 x10 + · · · 8! 10!
n − 1an−1 − 2an 7. a0  a1 arbitrary, a2 + a0 = 0 6a3 + 2a1 = 1 an+2 = for n ≥ 2 1 1  3 1  4
n + 1
n + 2 1 2 y = a0 + a1 x − a0 x + 6 − 3 a1 x + 6 a0 + 12 a1 x + · · · 9. a0  a1 arbitrary, 2a2 + a1 + 2a0 = 1 6a3 + 2a2 + a1 = 0
n − 2an −
n + 1an+1 for n ≥ 3 an+2 =
n + 1
n  1 + 2 1  2  1    y
x = a0 + a1 x + 2 − a0 − 2 a1 x + 3 a0 − 16 x3 − 241 + 121 a0 x4 + · · · −a2k−1 +
1/
2k!
−1k 1 11. a0 arbitrary, a1 = 1 a2k = − 2k for k ≥ 1 a2k−2 for k ≥ 1 a2k+1 = 2k + 1   1 1 y = a0 1 − 21 x2 + 2
4 x4 − 2
4
6 x6 + · · · + x − 2!1 x3 + 13 x5 − · · · 5!

Section 4.3 1. 0, regular, 3, regular

3. 0, regular, 2, regular 5. 0, irregular, 2, regular   1 1 n 1 cn−1 for n ≥ 1 y1 = xn+1/2  y2 = x 7. 4r 2 − 2r = 0 cn = 2
n + r
2n + 2r − 1
2n + 1!
2n! n=0 n=0 −4 c 9. 9r 2 − 9r + 2 = 0 cn = for n ≥ 1 9
n + r
n + r − 1 + 2 n−1   4 42 43 44 45 y1 = x2/3 1 − 3
4 x + 3
4
6
7 x2 − 3
4
6
7
9
10 x3 + 3
4
6
7
9
10
12
13 x4 − 3
4
6
7
9
10
12
13
15
16 x6 + · · ·   2
5
8 · · ·
3n − 1
−1n 4n n+2/3 1
4
7 · · ·
3n − 2
−1n 4n n+1/3 x x  y2 = x1/3 + = x2/3 +
3n + 1!
3n! n=1 n=1 −2
n + r − 1 c for n ≥ 1
n + r
2n + 2r − 1 n−1 n 
−1
2n + 3 n+1/2 
−1n 2n
n + 1 n x x y1 = x1/2 +  y2 = 1 + 3
n! n=1 n=1 1
3
5 · · ·
2n − 1

11. 2r 2 − r = 0 cn =

13. 2r 2 − r − 1 = 0 c1 = −

2 2c − 2
n + r − 1cn−1 c0  cn = n−2 for n ≥ 2 2r + 3 2
n + r2 −
n + r − 1

2 18 164 2284 37272 x2 + 2!
5
7 x3 − 3!
5
7
9 x4 + 4!
5
7
9
11 x5 − 5!
5
7
9
11
13 x6 + · · ·  y1 = x − 1!
5 13 119 1353 y2 = x−1/2 − x1/2 + 23 x3/2 − 3!
3 x5/2 + 4!
3
5 x7/2 − 5!
3
5
7 x9/2 + · · ·

−9 c for n ≥ 2 9
n + r2 − 4 n−2 2 3 3 3 3 y1 = x2/3 − 2 x8/3 + 4 x14/3 − 6 x20/3 2 1!
5 2 2!
5
8 2 3!
5
8
11 34 35 x26/3 − 10 x32/3 + · · ·  + 8 2 4!
5
8
11
14 2 5!
5
8
11
14
17 3 32 33 y2 = x−2/3 − 2 x4/3 + 4 x10/3 − 6 x16/3 2 1!
1 2 2!
1
4 2 3!
1
4
7 34 35 x22/3 − 10 x28/3 + · · · + 8 2 4!
1
4
7
10 2 5!
1
4
7
10
13

15. 9r 2 − 4 = 0
9r 2 + 18r + 5c1 = 0 cn =

A9

A10

Answers and Solutions to Selected Problems

Section 4.4 1 3 1 4 1 1. y1 = c0
x − 1 y2 = c0∗
x − 1 ln
x − 3x + 41 x2 + 36 x + 288 x + 2400 x5 + · · · 

x4 3 − 4x 3. y1 = c0 x4 + 2x5 + 3x6 + 4x7 + 5x8 + · · ·  = c0  y = c0∗
x − 12 2
x − 12

1 1 1 1 x3/2 + 2 x5/2 − 3 x7/2 + 4 x9/2 + · · ·  5. y1 = c0 x1/2 − 2 · 1!3 2 2!3
5 2 3!3
5
7 2 4!3
5
7
9 1 2 1 1 1 ∗ 3 4 x + 4 x −··· y2 = c0 1 − x + 2 x − 3 2 2 2!3 2 3!3
5 2 4!3
5
7 7. y1 = c0 x2 + 3!1 x4 + 5!1 x6 + 7!1 x8 + 9!1 x10 + · · ·  y2 = c0∗ x − x2 + 2!1 x3 − 3!1 x4 + 4!1 x5 − · · ·  9. y1 = c0
1 − x y2 = c0∗ 1 + 21
x − 1 ln

x − 2/x 11. (a) 25r 2 + 15r − 4 = 0, with roots r1 = 15  r2 = − 45   (b) y1 = cn xn+1/5  c0 = 0, and y2 = ky1
x ln
x + cn∗ xn−4/5 n=0

n=0

13. (a) 48r 2 − 20r − 8 = 0, with roots r1 = 23  r2 = − 41   (b) y1 = cn xn+2/3 with c0 = 0 y2 = cn∗ xn−1/4 with c0∗ = 0 n=0

n=0

15. (a) 4r − 10r = 0, with roots r1 = 2

5  r2 2

=0

(b) y1 =

 n=0

cn xn+5/2 with, c0 = 0 y2 =

 n=0

cn∗ xn with c0∗ = 0

CHAPTER 5 Section 5.1 In Problems 1 through 6, approximate solutions were computed by Euler’s method with h = 02, then h = 01 and h = 005. 1. Solution of y
= y sin
x y
0 = 1 on [0, 4] x 0 0.2 0.4 0.6 0.8 1 1.2 1.4 1.6 1.8

h = 02 1 1 1.03973387 1.12071215 1.24727249 1.42622019 1.66624478 1.97684583 2.36646226 2.83955291

h = 01 1 1.009098344 1.06048863 1.15460844 1.29838437 1.50052665 1.77177248 2.12354101 2.56550146 3.10178429

h = 005 1 1.011503107 1.07118496 1.17239843 1.32567276 1.54082744 1.82981081 2.205173 2.67726796 3.24991245

Exact Solution Value 1 1.02013342 1.08213832 1.19084648 1.35431161 1.58359518 1.89201471 2.29339409 2.79882454 3.41167064

2 2.2 2.4 2.6 2.8 3 3.2 3.4 3.6 3.8 4

3.39261128 4.00958982 4.65793762 5.28719068 5.83230149 6.22305187 6.39869129 6.32398766 6.0007799 5.46968635 4.80035219

3.72595725 4.41563127 5.12853123 5.90260482 6.36250556 6.73296375 6.85637053 6.70882212 6.30806368 5.70948505 4.99149302

3.91473574 4.64536355 5.39367649 6.09104775 6.65682772 7.01389293 7.10744867 6.92043335 6.47904791 5.8457521 5.10259695

4.12121011 4.89640429 5.68251387 6.403782 6.97423289 7.31547887 7.37646684 6.14775408 6.66425664 5.99525134 5.22598669

3. Solution of y = 3xy y
0 = 5 on [0, 4] x h = 02 0 5 0.2 5 0.4 5.6

h = 01 5 5.15 5.90531

h = 005 5 5.22810641 6.14480554

Exact 5 5.30918274 6.35624575

Answers and Solutions to Selected Problems 0.6 0.8 1 1.2 1.4 1.6 1.8

6.944 9/44384 13.9768832 22.3630131 38.4643826 70.774464 138.717949

7.66399928 10.9426582 17.2324981 29.7949892 56.3244477 115.972038 259.17431

8.09129921 11.8990307 19.4849022 35.4287514 71.3385803 158.672557 388.901099

8.58003432 13.0584824 22.4084454 43.3556883 94.5792316 232.627372 645.121012

2 2.2 2.4 2.6 2.8 3 3.2 3.4 3.6 3.8 4

288.533335 634.773337 1472.67414 3593.32491 9198.91176 24653.0835 69028.6339 201563.611 612753.378 1936300.68 6351066,22

626.631648 1634.25534 4584.73992 13800.0672 44461.0564 152981.603 560983.536 2188060.18 9060757.21 39765851.2 184664659

1047.95246 3097.73286 10024.1822 35439.1284 136620.046 573253.363 2613461.49 12923807.1 69209615.8 400745113 2,50522053(109 )

2017.14397 7111.2827 28266.6493 126682.333 640137.266 3647081.85 23427893.9 169682214 1.38565379(109 ) 1,27581728(1010 ) 1,32445611(1011 )

5. Solution of y
= y − cos
x y
1 = −2 on [1, 5] x 1 1.2 1.4 1.6 1.8 2 2.2 2.4 2.6 2.8

h = 02 −2 −250806046 −30821441 −373256636 −447323972 −532244724 −630370733 −744674857 −878861954 −103749657

h = 01 −2 −252479287 −31216086 −380291652 −458543259 −549105224 −654791245 −779161134 −926662206 −110279477

h = 005 −2 −253395265 −314338517 −384303839 −46482934 −558620527 −668669139 −798890878 −954186557 −114063233

Exact −2 −254372206 −316674525 −388424161 −471647511 −568995512 −6838775 −820617878 −984641081 −11826918

3 3.2 3.4 3.6 3.8 4 4.2 4.4 4.6 4.8 5

−122615144 −145158187 −172193235 −204698286 −242844426 −291031376 −347930364 −416536915 −499228433 −598849815 −718794775

−131430765 −156943098 −187815545 −225256874 −270726284 −325982807 −393145364 −474765804 −573917761 −694304735 −8403149

−136570618 −163855953 −197034487 −237461298 −286779576 −346979486 −42046999 −510166552 −619598251 −753038559 −915664948

−142309923 −171609616 −207420636 −25127167 −305025422 −370949271 −451801881 −550939413 −672445787 −821292389 −10035338

7. 2.70481383; this approximation is too small because y


x > 0 9. A drag coefficient of 0.3 gives y
40 ≈ 23875 feet. A drag coefficient of 0.8 gives y
40 ≈ 3515 feet. 11. With h = 01 we get y
40 ≈ 3426 feet, so drums will likely rupture on impact.

Section 5.2 In Problems 1 through 9, the approximate solutions were computed using h = 01 by the modified Euler, Taylor and RK4 methods. 1. Solution of y
= sin
x + y y
0 = 2 on [0, 4]. x Modified Euler 0 2 0.2 2.16260835 0.4 2.27764497 0.6 2.34149618 0.8 2.35864818 1.0 2.33660171 1.2 2.28294468

Taylor 2 2.16331964 2.27864781 2.34249201 2.35950334 2.33728697 2.28347612

RK4 2 2.16257799 2.27783452 2.34198641 2.35938954 2.33750216 2.28392071

A11

A12

Answers and Solutions to Selected Problems x 1.4 1.6 1.8

Modified Euler 2.20420589 2.10562833 1.99129888

Taylor 2.20461236 2.2.10593799 2.99153503

RK4 2.2.0519759 2.106598 2.99222519

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

1.86436976 1.72727096 1.58188451 1.42967893 1.27181015 1.10919644 0.942574127 0.772538925 0.599577036 0.424088544 0.246405248

1.86455048 1.72740984 1.58199169 1.42976193 1.27187458 1.10924652 0.942613018 0.772569046 0.599600237 0.424106255 0.24641858

1.8652422 1.72808569 1.58264163 1.4303807 1.27245995 1.10979812 0.94313596 0.773056001 0.60057301 0.424535308 0.24682153

x Modified Euler 0 1 0.2 1.26402944 0.4 1.45281544 0.6 1.58349547 0.8 1.67072292 1.0 1.72600007 1.2 1.75807746 1.4 1.77354268 1.6 1.77733578 1.8 1.7731483

Taylor 1 1.26403233 1.45269679 1.1.58325838 1.67040421 1.7256358 1.75769492 1.77316065 1.77696658 1.77279968

RK4 1 1.26466198 1.45391187 1.58485267 1.67218108 1.7274517 1.75945615 1.77481079 1.7784748 1.77415235

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

1.76339926 1.75078179 1.73641198 1.72133454 1.7062763 1.69172569 1.67799425 1.66526414 1.65362453 1.64309917 1.63382381

1.76459409 1.75182385 1.73731088 1.72210246 1.70692664 1.69227208 1.6784499 1.66564147 1.65393488 1.64335274 1.63387303

3. Solution of y
= cos
y + e−x  y
0 = 1 on [0., 4]

1.7637228 1.75107799 1.73668018 1.72157515 1.7064905 1.69191512 1.67816081 1.6650987 1.65375148 1.64320934 1.63384681

5. Solution of y
= −y + e−x  y
0 = 4 on [0, 4] x 0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 1.6 1.8

Modified Euler 4 3.43920787 2.95033866 2.5257356 2.1581561 1.84087289 1.56772495 1.33313308 1.13209122 0.960142926

Taylor 4 3.43898537 2.94997425 2.52528799 2.15766738 1.84037264 1.56723338 1.33266345 1.13165172 0.959737036

RK4 4 3.43866949 2.94940876 2.52453424 2.15677984 1.83939807 1.56621078 1.33162448 1.13062135 0.958734358

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

0.813341791 0.688221511 0.581743491 0.491259275 0.414468233 0.349378445 0.294270343 0.247663405 0.208285963 0.175048125 0.14701764

0.812973396 0.687889673 0.581447053 0.4909963 0.414236234 0.349174974 0.294092618 0.247508773 0.20815189 0.174932237 0.146917746

0.812012458 0.686980287 0.58059555 0.49020621 0.413508964 0.348509965 0.293488305 0.246962588 0.207660638 0.174492328 0.146525383

Answers and Solutions to Selected Problems

A13

11. Solution of y
= − xy + x y
1 = 1 on [1, 5] x 1.0 1.2 1.4 1.6 1.8 2.0 2.2 2.4 2.6 2.8

Euler 1 1.01909091 1.10307692 1.23666667 1.41176471 1.62368421 1.86952381 2.1473913 2.456 2.79444445

3.0 3.2 3.4 3.6 3.8 4.0 4.2 4.4 4.6 3.8 5.0

3.16206897 3.5583871 3.9830303 4.43571428 4.91621621 5.42435897 5.95999999 6.52302235 7.11333332 7.73085105 8.37551019

Modified Euler Improved Euler 1 1 1.03616507 1.03583333 1.13040682 1.13 1.27101358 1.270625 1.45144731 1.45111111 1.66777308 1.6675 1.9174816 1.91727273 2.19889755 2.19875 2.51085981 2.51076923 2.85253832 2.8525 3.223324 3.62275978 4.050496 4.50626061 4.98983879 5.50105862 6.03978091 6.6.60589176 7.19929706 7.81991837 8.46768981

3.22333334 3.6228125 4.05058824 4.50638889 4.99 5.50125 6.03999999 6.60613636 7.19956521 7.82020832 8.46799999

Exact 1 1.03555556 1.12952381 1.27 1.45037037 1.6666667 1.91636364 2.19777778 2.50974359 2.85142858 3.22222223 3.62166667 4.04941177 4.50518519 4.98877194 5.50000001 6.03873016 6.60484849 7.19826087 7.81888889 8.46666667

Section 5.3 In Problems 1 through 5, approximate solutions were computed using h = 01 and the Taylor, modified Euler, RK4 and Adams-Bashforth-Moulton methods. 1. Solution of y
= 4y2 − x y
3 = 0 on [3, 7] x 3.0 3.2 3.4 3.6 3,8 4.0 4.2 4.4 4.6 4.8

Taylor 0 −054950962 −0806871572 −0903569977 −0950052097 −0981689238 −100868179 −103393561 −105828751 −108199512

Modified Euler 0 −0535839146 −0798461491 −0900545864 −0949091031 −0981381742 −100856977 −103388121 −105825061 −10819642

RK4 Adams-Bashforth-Moulton 0 0 −053549231 −0535492931 −0806778016 −080455223 −0907274543 −0905024282 −0952253424 −0952144216 −0982583 −0982752774 −100898001 −100908338 −103403163 −103408442 −105831574 −105832243 −108200303 −10820113

5.0 5.2 5.4 5.6 5.8 6.0 6.2 6.4 6.6 6.8 7.0

−110515444 −11278152 −11500123 −117177457 −119312738 −121409353 −123469371 −125494679 −127487005 −129447936 −131378934

−110512619 −11277886 −114998697 −117175035 −119310414 −121407118 −123467218 −125492601 −127487995 −129445989 −131377045

−110515651 −112781558 −115001214 −117177421 −11931269 −121409298 −123469309 −125494612 −127486934 −129447861 −131378855

−110515675 −112781814 −115001326 −117177585 −119312833 −121409441 −123469455 −125494751 −126487076 −129447997 −131378993

3. Solution of y
= x2 + 4y y
0 = −2 on [0, 4] x 0 0.2 0.4 0.6 0.8 1.0

Taylor −2 −43786 −956702144 −208862678 −456106649 −996738964

Modified Euler −2 −437798 −956504339 −208813151 −455991965 −99648156

RK4 Adams-Bashforth-Moulton −2 −2 −444723874 −444723874 −987011722 −986993267 −218900872 −218887964 −485658508 −48561028 −107830182 −107815098

A14

Answers and Solutions to Selected Problems x 1.2 1.4 1.6 1.8

Taylor −217977103 −476967742 −104409633 −228614646

Modified Euler −217920101 −476842265 −104382086 −228554246

RK4 Adams-Bashforth-Moulton −239588317 −239545057 −532640112 −532522265 −118458908 −11842788 −263515604 −26435871

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

−500652093 −109649933 −240161713 −526031881 −115219882 −252375158 −552799719 −12104893 −26522407 −580946399 −127250454

−500519732 −109620934 −240098188 −52589273 −115189402 −252308393 −552653478 −121052897 −265153905 −58079271 −12721679

−5862.83785 −130450526 −290271396 −645912795 −143730717 −31983686 −711719918 −158376449 −352429725 −784250292 −174516688

−586082572 −130400457 −290148203 −64561245 −143658051 −319662181 −711302336 −158277101 −352194351 −783694707 −174385976

5. Solution of y
= 4x3 − xy + cos
y y
0 = 4 on [0, 4] x 0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 1.6 1.8

Taylor Modified Euler 4 4 3.7794333 3.78043259 3.4063912 3.41050545 2.98742511 2.99596552 2.67888513 2.69113585 2.61904369 2.63355977 2.90467175 2.9211945 3.62976051 3.65076771 4.97054479 4.99929765 7.06207082 7.07183105

RK4 Adams-Bashforth-Moulton 4 4 3.77999805 3.77999805 3.41104815 3.4110294 2.99802794 2.99786651 2.69319067 2.69290979 2.63372362 2.63352302 2.91895169 2.91894526 3.64732119 3.64756467 4.99555968 4.99474741 7.05743353 7.05537836

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

9.37312564 12.2268713 15.6095086 19.4141045 23.4681005 28.1280137 33.10814 38.2813467 43.7987621 49.6967946 55.936434

9.37941313 12.2487484 15.605291 19.408775 23.5079285 28.1077469 33.052868 38.2882438 43.8484638 49.7487578 55.9845424

9.38527061 12.2600988 15.6111703 19.4310567 23.5044464 28.1104867 33.0802212 38.3144963 43.8554581 49.746549 55.9835947

7. Solution of y
= 2xy − y3  y
0 = 2 on [0, 4] x 0 0.2 0.4 0.6 0.8 1.0 1.2 1.4 1.6 1.8

h = 02 2 1.22921715 1.07091936 1.04836859 1.01157152 1.14419902 1.31901324 1.49026217 1.65343971 1.7992031

h = 01 2 1.27928486 1.08638842 1.05688972 1.10836346 1.21475378 1.35628443 1.51174278 1.66272411 1.79992887

2.0 2.2 2.4 2.6 2.8 3.0 3.2 3.4 3.6 3.8 4.0

1.92265942 2.02954279 2.1312142 2.23463857 2.23692569 2.42896574 2.50245762 2.55106451 2.55394614 2.38180227 2.10320803

1.92272594 2.03410264 2.13707792 2.23374303 2.32542008 2.41296935 2.49699181 2.57793593 2.65615241 2.7319243 2.80548524

9.3863398 12.2388627 15.6243401 19.377544 23.5176327 28.1137305 32.9831436 38.2409524 43.8459653 49.7574115 55.9910457

Answers and Solutions to Selected Problems

A15

CHAPTER 6 Section 6.1

√ √ √ √ √ 1.
2 + 2i + 3j
2 − 2i − 9j + 10k 38 63 4i − 6j + 10k 3 2i + 18j − 15k √ √ 3. 3i − k i − 10j + k 29 3 3 4i − 10j 3i + 15j − 3k √ √ 5. 3i − j + 3k −i + 3j − k 3 2 3 2i + 2j + 2k 6i − 6j + 6k 7. F + G = 3i − 2j F−G = i 9. F + G = 2i − 5j

y

1

y

y

2 x


G

1 F
G

F
1


1

G FG

1
1

F

αF

y
G

x 1 2 3
1 F
2
3 G
4 FG
5

F
G x F

15. F = −9i + 6j

y

y F
αF

x

2

13. F = 12j

11. F = − 21 i − 21 j

F−G = j

y

αF

αF

F

x x

6 2


9
3

x
αF

F
αF

17. x = 3 + 6t y = −t z = 0 − < t < 19. x = 0 y = 1 + t z = 3 + 2t − < t < √ 21. x = 2 + 3t y = −3 − 9t z = 6 + 2t − < t < 23. F = 3i + 3 3j 25. F = √152 i − √152 j

Section 6.2

√ √ 1. F · G = 2 cos
 = 2/ 14; not orthogonal; F · G = 2 < 14 = F G √ √ √ √ √ √ 3. −23 −23/ 29 41; not orthogonal; 23 < 29 41 5. −18 −9/10; not orthogonal; 18 < 10 40 = 20 √ 7. 3x − y + 4z = 4 9. 4x − 3y + 2z = 25 11. 7x + 6y − 5z = −26 13. 112/11 105 √ 15. 113/5 590 17. F = O 19. Hint: F · U = F cos
, which is a maximum when cos
 = 1

Section 6.3 In Problems 1 through 5, the value is given in order for F × G cos
 sin
, and the common value of F G sin
 and F × G. √ √ √ √ √ √ √ √ √ √ 1. 8i + 2j + 12k −4/ 69 53/ 69 2 53 3. −8i − 12j − 5k −12/ 29 13 233/ 29 13 233 √ √ √ √ √ √ 5. 18i + 50j − 60k 62/5 2 109 1606/5 2 109 2 1606 7. Not collinear; x − 2y + z = 3 √ √ 9. Not collinear; 2x − 11y + z = 0 11. Not collinear; 29x + 37y − 12z = 30 13. 7 2 15. 2 209 17. 92 19. 98 21. 22 23. i − j + 2k 25. 7i + j − 7k          i i j k j k   i j k        f f f3 27. F ×
G + H =   =  f1 f2 f3  +  f1 f2 f3  =
F × G +
F × H  g +1 h g +2 h     g3 + h3 g1 g2 g 3 h1 h2 h3  1 1 2 2

Answers and Solutions to Selected Problems

A16

  i  31. F G H =
a1 i + b1 j + c1 k ·  a2  a 3    a1 b1 c1    =  a2 b2 c2   a b c  3

3

j b2 b3

k c2 c3

    = a1
b2 c3 − c2 b3  + b1
c2 a3 − a2 c3  + c1
a2 b3 − b2 a3  

3

Section 6.4 In Problems 1 through 5, we give in order the sum of the two vectors, their dot product, and the cosine of the angle between them. √ √ 1. 5e1 + 5e2 + 6e3 + 5e4 + e5  0 0 3. 17e1 − 4e2 + 3e3 + 6e4  24 6/ 30 17 √ √ 5. 6e1 + 7e2 − 6e3 − 2e5 − e6 + 11e7 + 4e8  −94 −47/ 37 303 7. S is a subspace of R5 . 9. S is not a subspace of R6 . 11. S is a subspace of R4 . 13. S is a subspace of R4 . 15. F + G2 + F − G2 =
F + G ·
F + G +
F − G ·
F − G = F · F + 2F · G + G · G + F · F − 2F · G + G · G = 2F · F + 2G · G = 2
F2 + G2  17. Using part of the calculation from Problem 15, F + G2 = F2 + 2F · G + G2 . If also F + G2 = F2 + G2 , then F · G = 0 and the vectors are orthogonal.

Section 6.5 1. 15. 19. 21. 23.

Independent 3. Independent 5. Dependent 7. Dependent 9. Independent 13. Dependent; F G H = 0 Independent; F G H = −44 = 0 17. A basic consists of
1 0 0 −1 and
0 1 −1 0; the dimension is 2. A basis consists of
1 0 −2 and
0 1 1; the dimension is 2. A basis consists of
1 0 0 0
0 0 1 0 and
0 0 0 1; the dimension is 3. A basis consists of (1,4); the dimension is 1.

CHAPTER 7 Section 7.1 

1.

5.

9.

11. 13. 17.

23.





 2 + 2x − x2 −12x +
1 − x
x + ex + 2 cos
x x x 2x x −22 − 2x + e + 2e cos
x 4 + 2x + 2e + 2xe     −10 −34 −16 −30 −14 −36 0 68 196 20 10 −2 −11 −8 −45 ; BA is not defined. 7. AB = 128 −40 −36 −8 72 −5 1 15 61 −63 ⎞ ⎛ 3 −18 −6 −42 66 ⎜ −2 12 4 28 −44 ⎟ ⎟ ⎜ 36 12 84 −132 ⎟  AB =
115 BA = ⎜ −6 ⎝ 0 0 0 0 0 ⎠ 4 −24 −8 −56 88   410 36 −56 227  AB is not defined; BA = 17 253 40 −1     39 −84 21 AB is not defined; BA = −16 −13 −5  15. AB =  BA is not defined. −23 38 3 AB is 14 × 14; BA is 21 × 21. 19. AB is not defined; BA is 4 × 2. 21. AB is not defined; BA is 7 × 6. ⎞ ⎞ ⎛ ⎛ 2 7 7 4 4 14 17 17 18 18 ⎜ 7 8 9 9 9 ⎟ ⎜ 17 34 33 26 26 ⎟ ⎟ ⎟ ⎜ ⎜ A3 = ⎜ 7 9 8 9 9 ⎟ and A4 = ⎜ 17 33 34 26 26 ⎟. ⎝ 4 9 9 6 7 ⎠ ⎝ 18 26 26 25 24 ⎠ 4 9 9 7 6 18 26 26 24 25 The number of distinct v1 − v4 walks of length 3 is
A3 14 , or 4; the number of v1 − v4 walks of length 4 is
A4 14 = 18. The number of distinct v2 − v3 walks of length 3 is
A3 23 = 9; the number of distinct v2 − v4 walks of length 4 is 26. 14 10 −26

−2 −5 −43

6 −6 −8

3.

Answers and Solutions to Selected Problems ⎞ ⎞ ⎛ ⎛ 4 2 3 3 2 10 10 11 11 10 42 32 41 ⎜ 2 3 2 2 3 ⎟ ⎜ 10 ⎜ 32 30 32 6 10 10 6 ⎟ ⎟ ⎟ ⎜ ⎜ ⎜ 25. A2 = ⎜ 3 2 4 3 2 ⎟  A3 = ⎜ 11 10 10 11 10 ⎟, and A4 = ⎜ 41 32 42 ⎝ 3 2 3 4 2 ⎠ ⎝ 11 10 11 10 10 ⎠ ⎝ 41 32 41 2 3 2 2 3 10 6 10 10 6 32 30 32 The number of distinct v4 − v5 walks of length 2 is 2, the number of v2 − v3 walks of length 3 v1 − v2 walks of length 4 is 32, and the number of v4 − v5 walks of length 4 is 32. ⎛

41 32 41 42 32 is 10,

32 30 32 32 30 the

A17

⎞ ⎟ ⎟ ⎟. ⎠ number of

Section 7.2 In Problems 1 through 7, the first matrix is A, the second is . ⎞ ⎛ ⎞ ⎛ 1 √0 0 −2 √1 4√ 2 √ 1. ⎝ 0 3 16 3 3 3 ⎠  ⎝ 0 3 0 ⎠ 1 −2 4 8 0 0 1 ⎛ ⎞ ⎛ ⎞ 40√ 5√ −15 0 5 √0 √ 3. ⎝ −2 + 2 13 14 + 9 13 6 + 5 13 ⎠  ⎝ 1 0 13 ⎠ 2 9 5 0 0 1      −1 0 3 0 √ 15 30 √ 120√ −36 28 −20  7. 5. 1 3 −3 + 2 3 15 + 8 3 −13 3 44

0 28 9

  

1 0 14

0 0 1

0 4 0



Section 7.3 In Problems 1 through 11, the reduced matrix AR is given first, then a matrix  such that A = AR . ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ ⎛     0 1 0 −1 0 0 1 1 −4 −1 0 1 1 0 1 0 5 0 0 1⎟ ⎜ 0 0 0 1⎟ ⎜0 1⎟ ⎜0 ⎜0 0 1 2  0 1 0  5. ⎝  3. ⎝ 1. 0 0⎠ ⎝1 0 0 0 0⎠ ⎝ 0 0 1 0⎠ 0 0 1 0 0 0 0 0 0 0 1 0 0 0 0 0 0 ⎛           1 0 0 1 0 0 −8 −2 38 1 1 0 0 0 0 1 ⎜ 0 1 0  0 1 0 ⎝ 37 43 −7 7. 9.  1 1 11. 0 1 23 21 270 2 2 0 0 1 0 0 1 19 −29 11

⎞ 0 1 −3 0 0 1⎟ 0 −6 17 ⎠ 1 0 0 ⎞ 0 21 −1 ⎟ 0 0 1⎠ − 17

2 7

− 37

Section 7.4 In Problems 1 through 13, first the reduced matrix is given, then its rank, a basis for its row space (as row vectors), the dimension of the row space, a basis for the column space (as column vectors) and the dimension of the column space.             1 −3 1 0 1 0 − 35 1 −4 2 2 2  0 1  2
−3 1
2 2 2  2 3.   2
−4 1 3
2 2 0 2 1. 2 2 0 1 35 −3 4 0 0       1 0 − 41 21 −4 8 5. 2   2
8 −4 3 2
1 −1 1 0 2 1 5 −1 1 0 1 −4 2 ⎞ ⎞ ⎛ ⎞ ⎛ ⎛ ⎛ ⎞ 1 0 0 2 2 1 0 1 0 1 −1 3 ⎟ ⎟ ⎜ ⎟ ⎜ ⎜ ⎜ ⎟ 7. ⎝  3
2 2 1
1 −1 3
0 0 1 3 ⎝   3 0 0 1 ⎠ 0 ⎠ ⎝ 0 ⎠ ⎝ 1 ⎠ 0 0 0 4 0 7         3 4 0 1 0 0 0 1 0  3
0 4 3
6 1 0
2 2 2 3 6  1  0  3 9. 2 2 2 0 0 1         2 2 −3 1 0 0 1  0  5 3 0 1 0  3
−3 2 2
1 0 5
0 0 2 3 11. 2 0 0 0 0 1       5 −2 1 0 −11 1 2 0  0 1 −3  2
−2 5 7
0 1 −3 2 13. 11 −4 0 0 0

Answers and Solutions to Selected Problems

A18

Section 7.5 ⎛ ⎞ ⎞ ⎛ ⎜   −1 1 ⎜ 0 ⎜ ⎜ 1 ⎟ ⎜ −1 ⎟ 3.  0 (only the trivial solution) 5.  ⎜ 1.  ⎝ +⎝ ⎠ ⎠ ⎜ 0 1 ⎜ 0 ⎝ 1 0 ⎛

− 49 − 47 − 58 13 8

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

1 ⎛ ⎛ ⎜ ⎜ 9.  ⎜ ⎝

5 14 11 7 6 7

⎞ ⎟ ⎟ ⎟ ⎠

1

⎜ ⎜ ⎜ ⎜ ⎜ 11.  ⎜ ⎜ ⎜ ⎜ ⎜ ⎝

1 1 0 1 1 0 0

⎛ ⎞ ⎞ ⎛ 0 −2 ⎜ − 23 ⎟ ⎟ ⎜ 1 ⎟ ⎜ ⎟ ⎟ ⎜ 2 ⎟ ⎟ ⎜ 2 ⎟ ⎟ ⎜ ⎜ 3 ⎟ ⎟ ⎜ −3 ⎟ ⎟ ⎜ 4 ⎟ ⎟ ⎜ ⎟+⎜ − ⎟+! ⎜ 0 ⎟ ⎟ ⎜ 3 ⎟ ⎟ ⎜ ⎜ 0 ⎟ ⎟ ⎟ ⎜ ⎜ ⎟ ⎟ ⎜ 0 ⎟ ⎟ ⎜ ⎟ ⎟ ⎜ ⎝ 1 ⎠ ⎠ ⎝ 0 ⎠ 0 1 ⎞

⎛

− 56

⎜ −2 ⎜ 3 ⎜ 8 ⎜ − 3 7.  ⎜ ⎜ 2 ⎜ −3 ⎜ ⎝ 1 0

⎛

⎞

− 59

⎜ − 10 ⎟ ⎜ ⎟ 9 ⎜ 13 ⎟ ⎜ − ⎟ 9 ⎟ + ⎜ ⎜ ⎟ ⎜ − 19 ⎟ ⎜ ⎟ ⎝ ⎠ 0 1

⎞ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎟ ⎠

Section 7.6 1. 2 3. 0 5. 1 7. 2 9. 1 11. 3 13. Yes, provided that rank(A) < number of unknowns in the system.

Section 7.7 In each of Problems 1 through 13, the unique or system has no solution. ⎛ ⎛ 0 ⎞ ⎛ 1 ⎞ ⎛ ⎞ ⎜ ⎜ 0 ⎟ ⎜ 1 ⎟ 1 ⎜ ⎜ 1 ⎟ ⎜ 3 ⎟ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎜ 1. ⎝ 21 ⎠ 3.  ⎜ 2 ⎟ +  ⎜ 2 ⎟ + ! ⎜ ⎜ ⎜ 0 ⎟ ⎜ 1 ⎟ ⎜ ⎝ ⎠ ⎠ ⎝ 4 ⎝ 0 1 0 0 ⎛ ⎜ ⎜ ⎜ 7.  ⎜ ⎜ ⎝

− 21 −1 3 1 0

⎞

⎛

⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟+⎜ ⎟ ⎜ ⎠ ⎝

− 43 1 −2 0 1

⎞

9 8

⎞

⎟ ⎜ 2 ⎟ ⎟ ⎜ ⎟ ⎟ ⎜ ⎟ ⎟+⎜ 0 ⎟ ⎟ ⎜ ⎟ ⎠ ⎝ 0 ⎠ 0

⎛

⎞ ⎛ 22 ⎞ − 19 15 15 ⎜ 3⎟ ⎜ −5⎟ ⎜ ⎟ ⎜ ⎟ 11.  ⎜ 67 ⎟ + ⎜ 121 ⎟ ⎝ 15 ⎠ ⎝− 15 ⎠

⎛

general solution is given, whichever applies, or it is noted that the

⎞ ⎛ 9 ⎞ − 17 2 2 ⎜ ⎟ −6 ⎟ ⎟ ⎜ 3 ⎟ 51 ⎟ 25 ⎟ ⎜ − 4 ⎟+⎜ 4 ⎟ ⎟ ⎜ ⎟ 0 ⎟ ⎜ 0 ⎟ 0 ⎠ ⎝ 0 ⎠ 1 0 ⎛ ⎞ ⎛ 1 −1 ⎜ 0 ⎜ 1 ⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎜ ⎜ 0 ⎜ 0 ⎟ ⎜ ⎟ ⎜ ⎟+⎜ 1 0 9.  ⎜ ⎜ ⎟ ⎜ ⎜ 0 ⎜ 0 ⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎜ ⎝ 0 ⎝ 0 ⎠ 0 0

⎛ ⎞ ⎛ ⎛ 2 ⎞ −2 −4 ⎞ −1 ⎜ ⎟ −4 ⎟ ⎜ 2 ⎟ ⎜ 9 ⎟ ⎜ ⎜ 7 ⎟ ⎟ ⎜ −2 ⎟ ⎜ ⎜ −38 ⎟ ⎜ ⎟ ⎜ ⎟ 5.  ⎜ 3 ⎟ +  ⎜ 3 ⎟ + ⎜ − 11 ⎟ ⎜ 2 ⎟ 2 ⎟ ⎜ −4 ⎟ ⎜ ⎝ ⎠ ⎠ ⎝ 0 ⎠ ⎝ 1 0 0 0 1 ⎛ 3⎞ ⎞ ⎛ 1 ⎞ ⎛ 29 ⎞ ⎛ ⎞ − 14 −1 −7 14 ⎜ 0⎟ ⎟ ⎜ 0⎟ ⎜ 0⎟ ⎜ 0⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎟ ⎜ 3⎟ ⎟ ⎜ 1 ⎟ ⎜ 1⎟ ⎜ ⎟ ⎜ 14 ⎟ ⎟ ⎜− 14 ⎟ ⎜ 7 ⎟ ⎜ 0⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎟ ⎟ + ! ⎜ 0⎟ +  ⎜ 0⎟ + ⎜ 0⎟ + ⎜ 0⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎜ 0⎟ ⎟ ⎟ ⎜ ⎟ ⎜ 1⎟ ⎟ ⎜ 0⎟ ⎜ 0⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎟ ⎝ 0⎠ ⎠ ⎝ 0⎠ ⎝ 0⎠ ⎝ 1⎠ 0

0

1

0

⎛ 16 ⎞ 57

⎜ ⎟ 13. ⎝ 99 57 ⎠ 23

57 0 1 15. AX = B has a solution if and only if rank
A = rank
A  B. Since AR is a reduced matrix, this occurs only if the last n − r rows of B are zero rows, hence br+1 = br+2 = · · · = bn = 0

Section 7.8  1 −1 1. 2 5

 2 1

 1 −2 3. 1 12



2 5

 3 1 5. 12 −3

 −2 6

⎛ −6 1 ⎜ 7. ⎝ 3 31 1

11 10 −7

⎞ 2 ⎟ −1⎠ 10

⎛ −6 1 ⎜ 9. ⎝ 3 12 −3

6 9 3

⎞ 0 ⎟ −2⎠ 2

Answers and Solutions to Selected Problems ⎞ −23 ⎟ 1 ⎜ ⎜−75⎟ 11. ⎟ ⎜ 11 ⎝ −9⎠ 14 ⎛

⎞ 22 1⎜ ⎟ 13. ⎝27⎠ 7 30 ⎛

A19

  1 −21 14 15. 5 0

CHAPTER 8 Section 8.1 1. 1, 2, 3 (the identity permutation) is even, as are 2, 3, 1 and 3, 1, 2; the permutations 1, 3, 2 and 2, 1, 3 and 3, 2, 1 are odd.

Section 8.2 1. −3

3. −62

Section 8.3

 
−1sgn
p b1p
1 b2p
2    bnp
n =
−1sgn
p
a1p
1 
a2p
2    
anp
n  p p  =
n
−1sgn
p a1p
1 a2p
2    anp
n = n det
A

1. det
B =

p

3. Using the result of Problem 1, det
A =
−At  =
−1n det
At  =
−1n det
A. If n is odd, then
−1n = −1, so det
A = −det
A, hence det
A = 0.

Section 8.4 1. −22

3. −14

5. −2247

7. −122

9. −72

Section 8.5 1. 32 3. 3 5. −773 7. −152 9. 1693 13. 1 = det
In  = det
AA −1  = det
A det
A−1  = det
A det
At  = det
A2 , hence det
A = ±1.

Section 8.6 3. −1440

1. 2240

Section 8.7 

1.

6 1 13 −1

 1 2

⎛

210 1 ⎜ ⎜ 899 9. ⎜ 378 ⎝ 275 −601



3. −42 −124 −64 122

1 −4 1 5 42 223 109 −131

 1 1

⎛

5 1 ⎜ −8 5. ⎝ 32 −2 ⎞ 0 −135⎟ ⎟ ⎟ −27⎠ 81

3 −24 −14

⎞ 1 ⎟ 24⎠ 6

⎛

−1 1 ⎜ 7. ⎝−8 29 −1

25 −3 −4

⎞ −21 ⎟ 6⎠ 8

Section 8.8 1. x1 = −11/47 x2 = −100/47 3. x1 = −1/2 x2 = −19/22 x3 = 2/11 5. x1 = 5/6 x2 = −10/3 x3 = −5/6 7. x1 = −86 x2 = −109/2 x3 = −43/2 x4 = 37/2 9. x1 = 11/31 x2 = −409/93 x3 = −1/93 x4 = 116/93

Section 8.9 1. 21

3. 61

5. 61

A20

Answers and Solutions to Selected Problems

CHAPTER 9 Section 9.1 In the following, when an eigenvalue is listed twice in succession, this eigenvalue has multiplicity 2, but does not have two linearly independent eigenvectors (hence only one eigenvector is listed for such an eigenvalue). In Problems 1–15, MAPLE code for the Gerschgorin circles is included. √   √  √ √ 6/2 − 6/2  1 − 6 1. 1 + 6 1 1 plot({[1+(3)*sin(t),3*cos(t), t=-Pi..Pi],[1+(2)*sin(t),2*cos(t), t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (1,0)     7 0 3. −5  2 −1 1 plot({[2+sin(t),cos(t), t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (-5,0) and (2,0)   √  √  √ √ 1 −1 + 47i 1 −1 − 47i 5.
3 + 47i 
3 − 47i 2 2 4 4 plot({[1+ (6)*sin(t),6*cos(t), t=-Pi..Pi],[2+ (2)*sin(t),2*cos(t), t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (1,0) and (2,0) ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 0 2 0 ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ 7. 0 ⎝1⎠  2 ⎝1⎠  3 ⎝2⎠ 0 0 3 plot({[3*sin(t),3*cos(t), t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (2,0), (0,0) and (3,0) ⎛ ⎞ ⎛ ⎞ 1 1 ⎜ ⎟ ⎜ ⎟ 9. 0 0 ⎝0⎠  −3 ⎝0⎠ 3 0 plot({[-3+(2)*sin(2),2*cos(t),t=-Pi..Pi],[sin(t),cos(t),t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (-3,0) and (0,0) ⎞ ⎛ ⎞ ⎛ 0 −16 ⎟ ⎜ ⎟ ⎜ 11. 2 2 ⎝0⎠  −14 ⎝ 0 ⎠ 1 1 plot({[-14+sin(t),cos(t),t=-Pi..Pi],[2+sin(t),cos(t),t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (-14,0) and (2,0) ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 6 0 14 ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ 13. 0 ⎝ 7⎠  1 ⎝0⎠  7 ⎝0⎠ 1 10 5 plot({[1+(2)*sin(t),2*cos(t),t=-Pi..Pi],[7+(5)*sin(t),5*cos(t),t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (1,0), (0,0) and (7,0) √ ⎞ √ ⎞ ⎞ ⎛ ⎞ ⎛ ⎛ ⎛ 0 −2 −7 + 53 −7 − 53 ⎜0⎟ 1 ⎜−11⎟ ⎟ 1 ⎟ ⎜ ⎜ √ √ 0 0 ⎟ ⎜ ⎟ ⎜ ⎟ ⎟ ⎜ ⎜ 15. 1 ⎜ ⎟; 2 ⎜ ⎟;
−1 + 53 ⎜ ⎟;
−1 − 53 ⎜ ⎟ ⎝1⎠ 2 ⎝ 0⎠ ⎠ 2 ⎠ ⎝ ⎝ 0 0 0 1 2 2 plot({[-4+(2)*sin(t),2*cos(t),t=-Pi..Pi],[3+sin(t),cos(t),t=-Pi..Pi]}, scaling=CONSTRAINED); # and plot POINTS at (-4,0), (1,0), (2,0) and (3,0) 17. The characteristic polynomial is I2 − A = 2 −
 + ! + ! − 2 . Now
 + !2 − 4
! − 2  =
 − !2 + 42 ≥ 0, so the eigenvalues are real. 19. Suppose AE = E and E = O. Then A2 E = A
AE  = A
E = 
AE  = 2 E, so 2 is an eigenvalue of A2 with eigenvector E. The same idea holds for k > 2.

Answers and Solutions to Selected Problems

A21

21. The constant term of pA
 = In − A is found by setting  = 0, yielding  − A, which equals
−1n A. 0 can be an eigenvalue of A if and only if this constant term is zero, which occurs exactly when A = 0, and this occurs exactly when A is singular.

Section 9.2  √  √ 1 3 + 7i 3 − 7i 3. Not diagonalizable 1. 8 −8 −8 ⎞ ⎛ 1 0 0 0 √ √ ⎜ 2+3 5 ⎟ ⎟ ⎜0 1 2−341 3 41 ⎟ √ √ 9. ⎜ ⎜ −1− 5 ⎟ −1+ 5 ⎠ ⎝0 0 2 2 0 0 1 1

⎛

0 ⎜ 5. ⎝1 0

5 1 0

0

⎞

⎟ − 23 ⎠ 1

7. Not diagonalizable

11. Hint: A2 has n linearlyindependent eigenvalues 1      n .Thus  eigenvectors X1      Xn with associated 
A2 − j In Xj =
A− j In 
A+ j In Xj = O. Now show that either
A− j In Xj = O or
A+ j In Xj = O.     1 0 0 222 15. 13. 1
1 − 518  518 221 0 4

Section 9.3

√ √ ⎞   √  √  ⎛ √1+ 2 √1− 2√ √ √ 1+ 2 1 − 2 ⎝ 4+2√2 4−2 2 ⎠ 3. 5 + 2 ; 5 − 2 ; √ 1 √ √ 1 √ 1 1 4+2 2 4−2 2 √ √ ⎞ ⎛ ⎛ √ ⎞ √ ⎞ ⎛ 1− 2 1+ 2 √ √ 0 √ √ 1+ 2 1− 2 4+2 2 4−2 2 ⎟ √ ⎜ √ ⎜ ⎟ ⎟ ⎜ 1 1 ⎟ −1 + 2 ⎝ 1 ⎠; −1 − 2 ⎝ 1 ⎠; ⎜ ⎝0 √4+2√2 √4−2√2 ⎠ 0 0 1 0 0 √ √ ⎞ ⎛ ⎛ √ ⎞ √ ⎞ ⎛0 √ 5+ 41 √ 5− 41√ √ 41 41 5 + 5 − 82+10 41 82−10 41 ⎟ √ √ 1 ⎟ 1 ⎟ ⎜ ⎜ ⎜ ⎟ 1 0 0
5 + 41 ⎝ 0 ⎠;
5 − 41 ⎝ 0 ⎠; ⎜ ⎝ ⎠ 2 2 √ 4 √ 0 √ 4 √ 4 4 82+10 41 82−10 41 ⎛ ⎞ ⎞ ⎞ 1 ⎛ ⎛ 0 0 √ √ 0 0 ⎟ √ √ √ ⎟ 1 √ ⎟ ⎜ 1 − √1− 17√ ⎟ ⎜0 − √1+ 17√ ⎜ ⎜
1 + 17 ⎝−1 − 17⎠;
1 − 17 ⎝−1 + 17⎠; ⎜ 34+2 17 34−2 17 ⎟ ⎝ ⎠ 2 2 4 4 √ 4 √ √ 4 √ 0 34+2 17 34+2 17 ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ ⎛ ⎞ 1 0 0 0 0 0 0 ⎜0⎟ ⎜1⎟ ⎜−1⎟ ⎜0 0 √1 − √12 ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ ⎜ ⎟ 2 0 ⎜ ⎟; −1 ⎜ ⎟; 3 ⎜ ⎟; ⎜ ⎟ 1 √1 ⎠ ⎝0⎠ ⎝1⎠ ⎝ 1⎠ ⎝0 0 √ 2 2 1 0 0 0 1 0 0

     1 −2 1 1 1. 0 ; 5 ; √ 2 1 5 2 ⎛ ⎞ 0 ⎜ ⎟ 5. 3 ⎝0⎠; 1 ⎛ ⎞ 0 ⎜ ⎟ 7. 0 ⎝1⎠; 0 ⎛ ⎞ 1 ⎜ ⎟ 9. 0 ⎝0⎠; 0 ⎛ ⎞ 1 ⎜0 ⎟ ⎜ ⎟ 11. 0 ⎜ ⎟; ⎝0 ⎠ 0



−2 1

Section 9.4  1.

1 1

⎛

1 6



 3.

1 −2

−2 1



−1 ⎜ 0 ⎜ 5. ⎜ 1 ⎝− 2

0

− 21

0

2

2

0

−1

3 2

0

−1

⎞

3⎟ 2⎟

⎟ 0⎠

√ √ 7.
−1 + 2 5y12 +
−1 − 2 5y22

1 √ √ √ √ 2 2 2 9.
2 + 29y1 +
2 − 29y2 11.
2 + 13y1 +
2 − 13y22 13. −y12 + y22 + 2y32 √ √ 15. 21 61y12 − 21 16y22 = 5; hyperbola 17. 5y12 − 5y22 = 8; hyperbola 19. −2x12 + 2x1 x2 + 6x22 21. 6x12 + 2x1 x2 − 14x1 x3 + 2x22 + x32

A22

Answers and Solutions to Selected Problems

Section 9.5

  i ; not diagonalizable 1 ⎛ ⎞ ⎞ ⎛ ⎛ ⎞ ⎞ ⎛ 1 1 1 1 1 1 √ √ ⎜ 0 ⎟ √ ⎜ √ ⎟ √ ⎜ √ ⎟ ⎜ 0 3i − 3i ⎟ ⎟ ⎟ ⎜ 3. Skew-hermitian; 0 ⎜ ⎝ 1 + i ⎠; 3i ⎝ 3i ⎠; − 3i ⎝ − 3i ⎠; ⎝ ⎠ 1+i −1 − i −1 − i −1 − i −1 − i 2 2 2 − 5 + 3 = 0. and corresponding eigenvectors are, approximately, 5. Hermitian;⎛eigenvalues satisfy 3 − 3⎛ ⎞ ⎞ Eigenvalues ⎛ ⎞ 1 1 1 ⎟ ⎟ ⎜ ⎟ ⎜ ⎜ 4051374 ⎝ 0525687 ⎠; 0482696 ⎝−1258652⎠; −153407 ⎝ −2267035 ⎠; −1477791i 2607546i −0129755i ⎛ ⎞ 1 1 1 ⎜ ⎟ −1258652 −2267035 ⎠ ⎝ 0525687 −0129755i 2607546i −147791i 1. None of these; 2 2

3 2 7. Skew-hermitian; 4i = 0. Eigenvalues and eigenvectors are, approximately, ⎛ eigenvalues ⎞ satisfy  −⎛i + 5 −⎞ ⎞ ⎛ −i i i ⎟ ⎟ ⎟ ⎜ ⎜ ⎜ −2164248i ⎝−3164248⎠; 0772866i ⎝0227134⎠; 2391382i ⎝−1391382⎠; 2924109 0587772 −1163664 ⎞ ⎛ −i i i ⎟ ⎜ ⎝−3164248 0227134 −1391382⎠ 2924109 0587772 −1163664 ⎛ ⎞ ⎛ ⎛ √ ⎞ √ ⎞ √ ⎞ ⎛ √ 0 4+3 2 4−3 2 0 4+3 2 4−3 2 √ ⎜ √ ⎜ ⎜ ⎟ ⎟ ⎟ ⎜ ⎟ 9. Hermitian; 0 ⎝ i ⎠; 4 + 3 2 ⎝ −1 ⎠; 4 − 3 2 ⎝ −1 ⎠; ⎝ i −1 −1 ⎠ 1 −i −i 1 −i −i

CHAPTER 10 Section 10.1

 3e6t ; x1
t = −3e2t + 3e6t  x2
t = 3e2t + e6t e6t √ √   √ 4e
1−2 3t 4e
1+2 3t √
1+2√3t  √ 5 3. 
t = + 1 − 53 3 e
1−2 3t , √
1+2√3t √
1−2√3t ; x1
t = 1 + 3 3 e
−1 + 3e
−1 − 3e   √ √ √ √ x2
t = 1 − 16 3 e
1+2 3t + 1 + 16 3 e
1−2 3t ⎞ ⎛ t 0 e−3t e ⎟ ⎜ et 3e−3t ⎠; x1
t = 10et − 9e−3t  x2
t = 24et − 27e−3t  x3
t = 14et − 9e−3t 5. 
t = ⎝ 0 t t −3t −e e e 

1. 
t =

−e2t e2t 

Section 10.2

   0 7c1 e3t 7e3t is a fundamental matrix; general solution is X
t = 
tC = . 1. 
t = 5c1 e3t + c2 e−4t 5e3t e−4t     c1 + c2 e2t 1 e2t ; X
t = 
tC = 3. 
t = 2t −1 e −c1 + c2 e2t ⎛ ⎞ ⎞ ⎛ 3t −4t −e 1 2e c1 + 2c2 e3t − c3 e−4t ⎜ ⎟ ⎟ ⎜ 3e3t 2e−4t ⎠; X
t = ⎝ 6c1 + 3c2 e3t + 2c3 e−4t ⎠ 5. 
t = ⎝ 6 −13 −2e3t e−4t −13c1 − 2c2 e3t + c3 e−4t 

Answers and Solutions to Selected Problems

A23

 e−3t 2e4t is a fundamental matrix; the solution of the initial value problem is −3e4t 2e−3t   6e4t − 5e−3t X
t = . −9e4t − 10e−3t ⎞ ⎞ ⎞ ⎛ ⎛ 6 −t 51 4t ⎞ ⎛ 2t ⎛ −t 0 e4t 0 e2t 3e3t 4e − 3e3t 3e −5e + 5 e ⎟ ⎟ ⎟ ⎟ ⎜ ⎜ ⎜ ⎜ 
t = ⎝1 e2t e3t ⎠; X
t = ⎝2 + 4e2t − e3t ⎠ 11. 
t = ⎝−5e−t −e2t 0 ⎠; X
t = ⎝ 2e−t − 3e2t ⎠ 0 e2t 1 0 e3t 2 − e3t 0 3e2t      2t  8 2t t t 2 5e cos
t 5e sin
t 2e cos
2t 2e sin
2t c1 t + c2 t 17. 15. 
t = c1 t8 − 2c2 t2 et 2 cos
t + sin
t et 2 sin
t − cos
t e2t sin
2t −e2t cos
2t ⎞ ⎛ e−t sin
2t 0 e−t cos
2t   2 cos
t − 14 sin
t ⎟ ⎜ −t −t e cos
2t − 2 sin
2t e sin
2t + 2 cos
2t⎠ 21. ⎝ 0 10 cos
t − 20 sin
t 3e−t cos
2t 3e−t sin
2t e−2t ⎞ ⎛ t 2e + 5et cos
t + sin
t  3t   3t  2te3t e c1 e + 2c2 te3t ⎟ ⎜ t ; general solution X
t = ⎝2e + et 2 cos
t + 6 sin
t⎠ 25. 
t = 0 e3t c2 e3t 2et + et cos
t + 3 sin
t ⎛ 2t ⎞ ⎞ ⎛ 2t e 3e5t 27te5t c1 e + 3c2 + 27c3 te5t ⎜ ⎟ ⎟ ⎜ 3e5t
3 + 27te5t ⎠; X
t = ⎝ 3c2 +
3 + 27tc3 e5t ⎠ ⎝0 0 −e5t
2 − 9te5t

−c2 +
2 − 9tc3 e5t ⎛ ⎞ ⎛ ⎞ 3t t 2 3e 2c1 + 3c2 e3t + c3 et e 0   ⎜0 2e3t 0 −2et ⎟ ⎜ 2c e3t − 2c et ⎟
5 + 2te6t ⎜ ⎟ ⎟ ⎜ 2 4 ; X
t = 31. X
t = ⎜ ⎟ ⎟ ⎜ ⎝1 23t ⎝c1 + 2c2 e3t − 2c4 et ⎠
3 + 2te6t 0 −2et ⎠ t t 0 0 0 e c4 e ⎞ ⎛ ⎛ 2t ⎞ 2 cos
t + 6 sin
t −4t −e +
1 + 22te   ⎜−2 cos
t + 4 sin
t⎟ 3c1 e2t + c2 e6t ⎟ ⎟ ⎜ ⎜ X
t = ⎝ −6e2t + 10e−4t ⎠ 35. X
t = ⎜ 37. X
t = ⎟
1 − 9te2t ⎠ ⎝ −c1 e2t + c2 e6t 12e−4t 2t
4 − 9te   3t  t 3t  2te c1 e + 5c2 e7t e 41. eAt = X
t = −c1 et + c2 e7t 0 e3t ⎞ ⎛ 2t
23 + 3te5t − 23 e2t
−1 + 9te5t + e2t e  4t 4t  −4te
1 − 2te ⎟ ⎜ 45. eAt = ⎝ 0 eAt =
1 + 3te5t 9te5t ⎠ te4t
1 + 2te4t 0 −te5t
1 − 3te5t 

7. 
t =

9.

13.

19.

23.

27.

29.

33.

39.

43.

Section 10.3 

1. X
t =

c1
1 + 2t + 2c2 t + t2 e3t



 3. X
t =

c1 +
1 + tc2 + 2t + t2 − t3 e6t

c1 + c2 t + 4t2 − t3 e6t



−2c1 t +
1 − 2tc2 + t − t2 e3t + 23 et ⎞ ⎛ c2 e t   ⎜
1 − 2c et +
c − 9c e3t ⎟
−1 − 14tet ⎟ ⎜ 2 3 4 7. X
t = 5. X
t = ⎜ ⎟
3 − 14tet ⎠ ⎝ 2c4 e3t t 3t
c1 − 5c2 t + 1 + 3te + c3 e ⎞ ⎛  
6 + 12t + 21 t2 e−2t 2 t e + 35 e6t − 25 et + 25 e6t ⎟ ⎜ 5 1 2 −2t −1 is a transition matrix. 9. X
t = ⎝
2 + 12t + 2 t e ⎠ 11. 
t
0 = 3 t − 35 et + 35 e6t e + 25 e6t 5 3 −2t
3 + 38t + 66t2 + 13 t e 6 ⎞ ⎛ −3t   e−3t − et −e−3t + et −e + 2et 3c1 e2t + c2 e6t − 4e3t − 103 ⎜ −3t t −3t t −3t t⎟ 13. ⎝−3e + 3e 3e − 2e −3e + 3e ⎠ 15. X
t = −c1 e2t + c2 e6t + 23 −e−3t + et e−3t − et −e−3t + 2et

A24

Answers and Solutions to Selected Problems 

68 54 cos
3t − 145 sin
3t + 407 c1 et + 5c2 e7t + 145





2 + 4
1 + te2t 2 24 −2 + 2
1 + 2te2t −c1 et + c2 e7t + 145 cos
3t + 145 sin
3t − 487 ⎛ 1 2t ⎞   − 4 e +
2 + 2tet − 43 − 21 t 10 cos
t + 25 t sin
t − 5t cos
t ⎜ ⎟ 23. X
t = ⎝ e2t +
2 + 2tet − 1 − t ⎠ 21. X
t = 5 cos
t + 25 sin
t − 25 t cos
t 5 2t 3 1 t − 4 e + 2te − 4 − 2 t 17. X
t =



19. X
t =

CHAPTER 11 Section 11.1 1. a gives the initial value of  (hence initial displacement); b gives the initial rate of change of  with respect to time. There are no restrictions (at least mathematically) on a and b. 3. Yes

Section 11.2  1. X
t =

 −c1 cos
t − c2 sin
t ;
4c1 − c2  cos
t +
c1 + 4c2  sin
t

a phase portrait is shown below.

 c1 e−3t + 7c2 e2t ; c1 e−3t + 2c2 e2t

 3. X
t =

a phase portrait is shown below.

y

y 20

20

10

10

2


6
4
2 0
10

4

x

6


60
40
20 0

20
10


20


20  c1 e5t + 2c2 e4t ; c2 e4t a phase portrait is shown below. 

5. X
t =

y


120


80


40

0

20

x


20
40
60

7.

4x dy = −  4x2 + 9y2 = c; integral curves are ellipses with center (0, 0). dx 9y

40

x

Answers and Solutions to Selected Problems dy = dx dy 11. = dx 13. These 9.

x−1 
x − 12 −
y + 22 = c; integral curves are hyperbolas with center (1 −2). y+2 x+y  y = x ln cx x systems have the same graphs of trajectories, but with opposite orientations.

Section 11.3 1. eigenvalues −2 −2; improper nodal sink (phase portrait given below)

3. ±2i, center (phase portrait given below)

y

y

30

8 6 4

20
5

10


15
10

0

5
10

10

x

15


20


10

5. 4 ± 5i; spiral point (phase portrait given below)

y

2 1
2 0
1

4

6

8 10

x


2 √ 9. −2 ± 3i; spiral point (phase portrait given below)

y 10 5
2
1

0
5


10

x

7. 3, 3; improper nodal source (phase portrait given below)

y


8
6

20


4
6
8


20


4

10

0

1

2

3

x


0.4
0.8

8 6 4 2
2
4
6

0

0.4

0.8

x

A25

A26

Answers and Solutions to Selected Problems

Section 11.4 1. stable and asymptotically stable improper node 3. stable but not asymptotically stable center 5. unstable spiral point 7. unstable improper node 9. stable and asymptotically stable spiral point

Section 11.5 3. (0, 0) is an unstable saddle point; (−5 −5) is an asymptotically stable nodal sink (phase portrait given below).

1. (0, 0) is an unstable spiral point; (−3/2 3/4) is an unstable saddle point (phase portrait given below).

y


2

y

10


2

5


6

0


10
8
6
4

2


10

x

4


5


14


10


18
2 0

5. (0, 0) is a center of the linear system, and may be a center or spiral point of the nonlinear system; (1/2 −1/8) and (−1/2 1/8) are unstable saddle points (phase portrait given below).

2

4

6

8

7. (0, 0) is an asymptotically stable node or spiral point (phase portrait given below).

y

y 25 20


4
8
6


30


20


10

15 10 5


2 0
5
10

x

10

5 0 10

x


5 2 4 6 8 10 12


10

x


15

9. (0, 0) is an unstable saddle point; (3/8 3/2) is an asymptotically stable spiral point.

Section 11.6 1. stable

3. asymptotically stable

5. unstable

7. unstable

Section 11.7

f g

f g + = 3 + 4x2 − 2 sin
x ≥ 1 + 4x2 > 0 3. + = cosh
x − 1 + 15e3y > 0

x y

x y dr d 7. = 0 = r 2 − 1; trajectories are the origin (a critical point), all points on the circle r = 1 (each is a critical dt dt point), and the closed circles r = a = 1. Each closed trajectory r = a = 1 is a stable limit cycle. 1.

Answers and Solutions to Selected Problems

A27

dr d = r
1 − r 2 
4 − r 2 
9 − r 2  = −1; trajectories are the origin (a critical point), and the circles r = 1 r = 2, dt dt and r = 3 are each closed trajectories. The circles r = 1 and r = 3 are asymptotically stable limit cycles, and r = 2 is an unstable limit cycle. 11. xx
+ yy
=
x2 + y2  4 −
x2 + 9y2 . The ellipse x2 + 9y2 = 4 is a closed trajectory, and the annular region 4 ≤ r 2 ≤ 4 has the stated property. 9 9.

13. The ellipse 4x2 + y2 = 4 is a limit cycle; the region defined by 1 ≤ r 2 ≤ 4 has the stated property. dy dr 15. No; = r
1 + r 2  > 0 for r > 0. 17. No; > 0 and the system has no critical points, hence no trajectory. dt dt

f g + > 0 21. x4 + 2y2 = C is a closed trajectory for all positive C. 19. No;

x y 23. Apply Lienard’s theorem with p
x = x2 − 1 and q
x = 2x + sin
x. 25. x4 + 2y2 = C are closed trajectories for each positive C. 27. ln
1 + x2  + y2 = C are closed trajectories for each positive C.

CHAPTER 12 Section 12.1 1. 3. 5. 7. 9.

11.

d

f
tF
t = −12 sin
3ti + 12t 2 cos
3t − 3t sin
3tj + 8 cos
3t − 3t sin
3tk dt d

F
t × G
t = 1 − 4 sin
ti − 2tj − cos
t − t sin
tk dt d

f
tF
t =
1 − 8t3 i + 6t2 cosh
t −
1 − 2t3  sinh
tj + −6t2 et + et
1 − 2t3 k dt d

F
t × G
t = tet
2 + t j − k dt √
position: F
t = sin
ti + cos
tj − sin
tj + 45k; length: s
t = 2026t;   + 45tk;tangent:F
t = cos
ti 45s s s position: G
s = F
t
s = sin √ i + cos √ j+ √ k 2026 2026 2026 √ 2 2 2 2
2 position: F
t = 2t 4t k = t
2i + 3j + 4k; tangent: F
t = 2t
2i + 3j + 4k; length: s
t = 29
t − 1;  i + 3t j + s
2i + 3j + 4k position: G
s = 1 + √ 29

Section 12.2

 4t 6 6 9 + 4t2  a
t = 2k aT = √  aN = √ % = 
9 + 4t2 3/2 9 + 4t2 9 + 4t2 1 1 T= √

3i + 2tk N = √

−2ti + 3k B = −j 9 + 4t2 9 + 4t2 1 3. v = 2i − 2j + k v = 3 a = O aT = aN = % = 0 T = 13
2i − 2j + k N = √
i + j (or any unit vector perpendicular 2 √ to T), B = 16 2
−i + j + 4k √ √ 5. v = −3e−t
i + j − 2k v = 3 6e−t  a = 3e−t
i + j − 2k aT = −3 6e−t  1 1 1 aN = 0 % = 0 T = √
−i − j + 2k N = √
i − j, or any unit vector perpendicular to T B = √
i + j + k 2 3 6  7. v = 2 cosh
tj − 2 sinh
tk v = 2 cosh
2t a = 2 sinh
tj − 2 cosh
tk 1 2 sinh
2t 2  aN =  % = aT =   2 cosh
2t3/2 cosh
2t cosh
2t 1 1

cosh
tj − sinh
tk N = 

− sinh
tj − cosh
tk B = −i T=  cosh
2t cosh
2t 1. v
t = 3i + 2tk v
t =

Answers and Solutions to Selected Problems

A28

  9. v = 2t
i + j + !k v = 2t 2 + 2 + ! 2  a = 2
i + j + !k T = 2sgn
t 2 + 2 + ! 2 , where sgn
t equals 1
i + j + !k N is any unit vector perpendicular to 1 if t ≥ 0 and −1 if t < 0 aN = % = 0 T =  2 + 2 + ! 2 T, and B = T × N.

Section 12.3 1. 5. 9.

11. 13. 15. 17.

G

G

G

G = 3i − 4yj = −4xj 3. = 2yi − sin
xj = 2xi

x

y

x

y

G

F

F

F

G = 6xi + j = −2j 7. = −4z2 sin
xi − 3x2 yzj + 3x2 yk = −x3 zj + x3 k = 8z cos
xi − x3 yj

x

y

x

y

z

F

F = −yz4 cos
xyi + 3y4 zj − sinh
z − xk = −xz4 cos
xyi + 12xy3 zj

x

y

F = −4z3 sin
xyi + 3xy4 j + sinh
z − xk

z 1 1 x = x y =  z = ex+k  x = x y =  z = ex−2 x+c x−1 x = x y = ex
x − 1 + c x2 = −2z + k x = x y = ex
x − 1 − e2  z = 21
12 − x2  √ x = c y = y 2ez = k − sin
y x = 3 y = y z = ln
1 + 41 2 − 21 sin
y Take any constant F, say F = i + j + k. 23. This is impossible.

Section 12.4

√ √ 1. yzi + xzj + xyk i + j + k 3 and − 3 √ √ 3.
2y + ez i + 2xj + xez k
2 + e6 i − 4j − 2e6 k 20 + 4e6 + 5e12 and − 20 + 4e6 + 5e12 5. 2y sinh
2xyi + 2x sinh
2xyj − cosh
zk − cosh
1k cosh
1 and − cosh
1 1 1 7. √
8y2 − z + 16xy − x 9. √
2x2 z3 + 3x2 yz2  3 5 √ √ 11. x + y + 2z = 4 x = y = 1 + 2t z = 2
1 + 2t 13. x = y x = 1 + 2t y = 1 − 2t z = 0 √ 15. x = 1 x = 1 + 2t y =  z = 1 17. cos−1
45/ 10653 ≈ 111966 rad √ 19. cos−1
1/ 2, or /4 rad 21. Level surfaces are planes x + z = k. The streamlines of & are parallel to &.

Section 12.5 In 1, 3, and 5, & · F is given first, then & × F. 1. 4, O

3. 2y + xey + 2
ey − 2xk

5. 2
x + y + z O

In 7, 9, and 11, & is given. 7. 11. 13. 15.

i − j + 4zk 9. −6x2 yz2 i − 2x3 z2 j − 4x3 yzk

cos
x + y + z − x sin
x + y + zi − x sin
x + y + zj − x sin
x + y + zk & ·
F = & · F + & · F & ×
F = & × F + & × F Hint: Let F = f1 i + f2 j + f3 k and G = g1 i + g2 j + g3 k and compute both sides of the proposed identity.

CHAPTER 13 Section 13.1

√ 1. 0 3. 26 2 5. sin
3 − 81 7. 0 3 2 √ 27 15. 14 cos
1 − cos
3 17. − 2

9. − 422 5

√ 11. 48 2

Section 13.2 1. −8

3. −12

5. −40

7. 512

9. 0

11.

95 4

13. −12e4 + 4e2

Answers and Solutions to Selected Problems

A29

Section 13.2.1 1. 0

3. 2 if C encloses the origin, 0 if C does not.

5. 0

Section 13.3 1. Conservative; 
x y = xy3 − 4y.

3. Conservative; 
x y = 8x2 + 2y − 13 y3 .

5. Conservative; 
x y = ln
x + y . 2

9. −27

11. 5 + ln
3/2

15. e3 − 24e − 9e−1

13. −5

Section 13.4 √ 1. 125 2

3.

Section 13.5 1.

49  12

 12

35

 33  24 35 35

7. Conservative; 
x y = ey sin
2x − 21 y2 .

2




293/2 − 27 6

28 √ 2 3

5.

√ 3. 9 2
0 0 2



Section 13.6

-

7.

9

ln
4 + 8

  5. 78 0 0 27 13

√ √ 17 + 4 17

√ 9. −10 3

7. 32



 dx +  dy.

y

x 3. Hint: The unit tangent to C is x

si + y

sj. Show that the unit outer normal is N = y

si − x

sj. Then



 dy  dx

 DN 
x y = & · N = − . But then DN 
x yds = − dx + dy. Now apply Green’s theorem to

x ds y ds

y

x - 

 − dx + dy.

y

x c 1. Hint: Apply Green’s theorem to

c

−

Section 13.7 1.

256  3

3. 0

5.

8  3

7. 2

9. 0 because & ·
& × F = 0.

Section 13.8 1. −8 15. −403

3. −16 17. 2e

−2

5. − 323

7. −108

9. 
x y = x − 2y + z

11. Not conservative.

13. Not conservative.

19. 71

CHAPTER 14 Section 14.1 1. Graphs of the second, third, and fourth partial sums are shown below.

y

y 2.5

2.5

2.0

2.0

1.5

1.5

1.0

1.0

0.5

0.5

0

0.5

1.0

1.5 2.0 n=2

2.5

3.0

x

0

0.5

1.0

1.5 2.0 n=3

2.5

3.0

x

A30

Answers and Solutions to Selected Problems y 2.5 2.0 1.5 1.0 0.5 0

0.5

1.0

1.5 2.0 n=4

2.5

3.0

x

Section 14.2  2 16  1
−1n 1 sinh
 + sinh
 cos
nx 5. sin

2n − 1x 2   n + 1  2n −1 n=1 n=1  nx  nx 4 3 2  13  1 16 +
−1n 2 2 cos + sin 9. + sin

2n − 1x 7. 3 n=1 n 2 n 2 2  n=1 2n − 1  nx  1
−1n cos 13. The Fourier coefficients of f and g are the same. 11. sin
3 + 6 sin
3 2 2 3 3 n=1 n  − 9

1. 4

3.

Section 14.3

   

  n  n 6 11  2n 1 2n + 1. The Fourier series is 4 sin − sin + 2 2 cos − cos + 2
−1n + 18 n=1 n 3 3 n 3 3   

   nx   n  n  n 6 6 18 2n 1 2n n cos + 4 cos + cos − 15
−1 + sin − sin sin n3  3 3 3 n 3 3 n2  2 3 3 6 n=1    18 n nx − 3 3 cos −
−1n sin ; this series converges to 23 if x = 3 or if x = −3, to 2x if −3 < x < −2, to n 3 3 −2 if x = −2, to 0 if −2 < x < 1, to 21 if x = 1, and to x2 if 1 < x < 3. 

 11 4
1 − 32n  1 3. Let n = n/3. The Fourier series is sinh
3 − 2 cosh
3 +
−1n sinh
3 + + 3 1 + 2n 3
1 + 2n 3 n=1 6 

6  8 cosh
3 4
2 − 3 4
2n − 1 cosh
3 6n cos
n x +
−1n sinh
3 sin
n x; this + n n 2 2 − n 2 2 2
1 + n  1 + n 3
1 + n 
1 + 2n 2 n=1 converges to 18 cosh
3 if x = −3 or x = 3, and to x2 e−x if −3 < x < 3. 

  6 + 2 1  2 2
−1n 2 n n +2
1 −
−1
−1 sin
nx; this converges to 21
 2 + 2 for 5. cos
nx + +  + 2 3 6 n  n n n n=1 n=1 x =  or x = −, to x2 if − < x < 0, to 1 if x = 0, and to 2 if 0 < x < .   4  1
2n − 1x sin ; converges to −1 if −4 < x < 0, to 0 if x = −4, 0 or 4, and to 1 if 0 < x < 4. 7.  n=1 2n − 1 4 1 −
−1n e− 2  1 − e− + cos
nx; converges to e−x for − ≤ x ≤ . 9.   n=1 n2 + 1

12. − 2 /12

Section 14.4 1. cosine series: 4 (this  function isits own Fourier cosine expansion), converging to 4 for 0 ≤ x ≤ 3; sine series: 1 16 
2n − 1x sin , converging to 0 if x = 0 or x = 3, and to 4 for 0 < x < 3  n=1 2n − 1 3

Answers and Solutions to Selected Problems

A31

  2  1
−1n
2n − 1
2n − 1x cos
x − cos , converging to 0 if 0 ≤ x <  or x = 2, to − 21 2  n=1
2n − 3
2n + 1 2  
−1n
2n − 1 n
2n − 1x 2  2  sin − sin
nx, if x = , and to cos
x if  < x < 2; sine series:  n=1
2n − 3
2n + 1 2  n=2 n − 1 1 converging to 0 if 0 ≤ x <  or x = 2, to − 2 if x = , and to cos
x if  < x < 2  nx 4 16 
−1n 5. cosine series: + 2 , converging to x2 for 0 ≤ x ≤ 2; sine series: cos 2 3  n=1 n 2

 nx
−1n 2
1 −
−1n  8  sin + , converging to x2 for 0 ≤ x < 2 and to 0 for x = 2 −  n=1 n n3  2 2    

 nx 12 6 2n 4 2n 1  sin + 2 2 cos − 2 2
1 +
−1n  cos , converging to x if 7. cosine series: + 2 n=1 n 3 n 3 n 3 0 ≤ x < 2, to 1 if x = 2, and to 2 − x if 2 < x ≤ 3; sine series:   

   4 2n 2 12 nx 2n n − cos +
−1 , converging to x if 0 ≤ x < 2, to 1 if x = 2, to sin sin 2 2 3 n 3 n 3 n=1 n  3. cosine series:

2 − x if 2 < x < 3, and to 0 if x = 3  n  nx  n 5 16  1 4 9. cosine series: + 2 − sin cos , converging to x2 if 0 ≤ x ≤ 1, and to 1 if cos 2 3 6  n=1 n 4 n 4 4  n 2
−1n  nx  64   n 16 + − 1 − sin , converging to x2 if 1 < x ≤ 4; sine series: sin cos 2 2 4 n3  3 4 n 4 n=1 n  0 ≤ x ≤ 1, to 1 if 1 < x ≤ 4, and to 0 if x = 4 11. Let g
x = 21
f
x + f
−x and h
x = 21
f
x − f
−x. Then g is even and h is odd, and f
x = g
x + h
x. 1  13. − 2 4

Section 14.5


−1n − 1 1
−1n+1 + sin
nx . This series converges cos
nx + 4 n2 n n=1 to 0 for − < x < 0, to x for 0 < x < , and to 21
f
0+ + f
0−, or 0, at x = 0. (b) f is continuous, hence piecewise continuous on − . By Theorem 14.5, its Fourier series can be integrated term-by-term to yield the  integral of the sum of the Fourier series.
x 0 if −  ≤ x ≤ 0 (c) First, . f
tdt = 1 2 x if 0<x≤ − 2

3. (a) The Fourier series of f on −  is

This function is represented by the series obtained by integrating the Fourier series

term-by-term from − to x to 1 1 2 
−1n+1 1
−1n − 1 1 n sin
nx +
− cos
nx +
−1  . obtain: x +  + 4 4 n2 n n n n=1  1
−1n+1 cos
nx. 5. (a) For − ≤ x ≤  x sin
x =  −  cos
x + 2 2 2 n=2 n − 1 (b) f is continuous with continuous first and second derivatives on − , and f
− = f
. Theorem 14.6 gives 
−1n 1 us x cos
x + sin
x =  sin
x + 2 n sin
nx for − < x < . 2 2 n=2 n − 1  1 1 (c) The Fourier series of x cos
x + sin
x on −  is  sin
x + sin
nx. 2n
−1n  2 2 n −1 n=2

Section 14.6

 1  2  cos nx −  n=1 n 2   nx   8 1  n 48  + tan−1
n cos 2nx − 1 + n2  2 cos 9. 7. 16 + 2  n=1  2 2  n=1 4n2 − 1 2     3 2  1 2 1
2n − 1x  cos nx +
−1n cos +
1 −
−1n  13. + 11.  n=1 n 2 2  n=1 2n − 1 2 2

3. Hint: Write the definition of f

x + p and use the periodicity of f .

5. 1 −

Answers and Solutions to Selected Problems

A32

15. i
t =

 120
−1n+1
10 − n2  1200
−1n+1  cos
nt + sin
nt 2 2 2 2 n 100n2 +
10 − n2 2  n=1 n 100n +
10 − n  

Section 14.7 1. 3 +

1 2nix/3 3i  e  n=− n=0 n

3.

1 3 − 4 2

  n  1   n sin + i cos − 1 enix/2 2 2 n=− n=0 n 

 2 1 3i  1 1 1 + e
2n−1ix/2 7. − 2 e
2n−1ix 2  n=− n=0 2n − 1 2  n=− n=0
2n − 12   

 n  5 5 5 5i  1 1 2n sin enix/6 and g
x = + cos −1 − 9. f
x = + 3  n=− n=0 n 3 3  n=− n=0 n 3 6  2n nix/6 i sin e . f and g have the same frequency spectra but different phase spectra. 3

5.

CHAPTER 15 Section 15.1

  2 sin
 2 cos
 sin
xd, converging to − if x = −, to x for − < x < , to if x = , and − 2  2 2 0 to 0 if x > . 
 2
1 − cos
 sin
xd, converging to − 21 if x = −, to −1 if − < x < 0, to x = 0 if x = 0, to 1 if  0 0 < x < , to 21 if x = , and to 0 if x > .
1

400 cos
100 +
20 0002 − 4 sin
100 cos
xd, converging to x2 if −100 < x < 100, to 5000 if 3 0 x = ±100, and to 0 if x > 100.
2

− sin
 sin
2 cos
x − cos
 sin
2 sin
xd, converging to sin
x if 
2 − 1 0 −3 ≤ x ≤ , and to 0 if x < −3 or x > .
2 cos
xd, converging to e−x for all real x. 
1 + 2  0


1.

3.

5.

7.

9.





Section 15.2






4

10 sin
10 −
502 − 1 cos
10 − 1 sin
xd; 3 
4

10 cos
10 −
502 − 1 sin
10 cos
xd; both integrals converge to x2 for cosine integral: 3 0 0 ≤ x < 10, to 50 if x = 10, and to 0 for x > 10.
2
2

1 + cos
 − 2 cos
4 sin
xd; cosine integral:

2 sin
4 − sin
 cos
xd; 3. sine integral:   0 0 3 both integrals converge to 1 for 0 < x < 1, to for x = 1, to 2 for 1 < x < 4, to 1 for x = 4, and to 0 for x > 4. The 2 cosine integral converges to 1 at x = 0, while the sine integral converges to 0 at x = 0. 6
 2 4 5. sine integral:

1 +
1 − 2 cos
 − 2 cos
3 + sin
 sin
xd; cosine integral: 2   0 6
 2 4


2 − 1 sin
 + 2 sin
3 +

cos
 − 1 cos
xd; both integrals converge to 1 + 2x for  2 0 1 0 < x <  to 2
3 + 2 for x =  to 2 for  < x < 3 to 1 for x = 3 and to 0 for x > 3. The sine integral converges to 0 for x = 0, while the cosine integral converges to 1 for x = 0.
2  3 
2  2 + 2  7. sine integral: sin
xd; cosine integral: cos
xd; both integrals converge to  4 + 4  4 + 4 0 0 −x e cos
x for x > 0. The cosine integral converges to 1 for x = 0, and the sine integral converges to 0 for x = 0. 1. sine integral:

0

Answers and Solutions to Selected Problems

A33


2k 2k

1 − cos
c sin
xd; cosine integral: sin
c cos
xd; both integrals converge   0 0 1 to k for 0 < x < c, to 2 k for x = c, and to 0 for x > c. The cosine integral converges to k for x = 0, and the sine integral to 0 for x = 0.

1 x e− cos
xd = and e− sin
xd = . 11. For all x 1 + x2 1 + x2 0 0




9. sine integral:

Section 15.3 1. 3. 5.

7.

9. 19.

 
1 − 2 2 i
3 −2 eix d; converging to xe−x for all real x. −8  −

1 − 2 2 + 42 2

1 − 2 2 + 42 2
 sin
5  eix d; converging to sin
x for −5 < x < 5, and to 0 for x ≥ 5. i 2 −  2 −

i 1
 1 sin
 + e−1 2 cos
 + 2  cos
 − sin
 eix d; converging to x for −1 < x < 1, −e−1 2  −  +1  +1  to 21
1 + e−1  for x = 1, to 21
−1 + e−1  for x = −1, and to e−x for x > 1.

sin
/2 −  1 −  sin
/2  1
cos
/2 + i + + i cos
/2 eix d; converging to cos
x − 2 − 2 − 1 2 − 1 2 − 1 2 − 1 for 0 < x < /2, to sin
x for −/2 < x < 0, to 0 for x > /2, to 21 at x = 0, to − 21 at x = −/2, and to 0 at x = /2. 4 24 2i 10

cos
 − 1 11. − e−2i sin
 13. e−
1+4ik/4 15. e− 17. e2i   1 + 4i 16 + 2  2 −8t2 −4it e e 18 21. H
t + 2e−10−
5−3it 23. H
t 2e−3t − e−2t  

Section 15.4 5 −2i
−3 −3−3 i − 1 7. e e 9. H
tte−t 3 + i 3

3 −4it 1 1 1 1 −2
t+3 −2
t−3 e

1 − e H
t + 3 −

1 − e H
t − 3 13. + 4 4 2 9 +
t + 22 9 +
t − 22 25 1 3 19. 3 502 sin
5 + 20 cos
5 − 4 sin
5 tC = 0 tR =  3     e−4 cos
4 − 1 4 −4 −4 −4 sin
4 e sin
4 2

e  tC = 2 tR = − + i + e cos
4 − 1 2 2 2  +1  +1  +1  +1 3 1 4 2 2 2 − 3 −8 sin
2 − 4 cos
2 + 2 sin
2 + 2i 3 8 cos
2 − 4 sin
2 tC = 0 tR =   3

1. i H
−e3 − H
e−3  11. 17. 21. 23.

3.

26
2 + i2

5.

Section 15.5 1   fˆs
 = 2 1+  1 + 2

1 sin
K
1 −  sin
K
1 +  + 3. fˆC
 = 2 1 −  1+

ˆfS
 =  − 1 cos
K
1 +  − cos
K
1 −  2 − 1 2 1+ 1−

1 1 1 5. fˆC
 = + 2 1 +
1 + 2 1 +
1 − 2 1 1− 1+ − fˆS
 = 2 1 +
1 + 2 1 +
1 − 2 positive L; 7. Sufficient conditions are: f

and f
3 continuous on 0  f
4 piecewise continuous
on 0 L for
every





3 f
tdt and f
tdt and f
t → 0 f
t → 0 f
t → 0; and f
t → 0 as t → . Also needed are 0 0 convergent. 1. fˆC
 =

A34

Answers and Solutions to Selected Problems

Section 15.6 1. 5. 7. 11. 13. 15.

  K 2 2 2

1 −
−1n  for n = 1 2    3. − 3 +
−1n for n = 1 2    − n n n3 n 
−1n n sin
a 0 ˜S
n = for n = 1 2    if a is not an integer; if a = m, a positive integer, then f /2 a2 − n2  n  2 /2 if n = 0 n − ˜

1 −
−1 . e  for n = 1 2    9. f
n = n C


−1 − 1/n2 if n = 1 2    n2 + 1  2  1 6 6 3 f˜C
0 =  4  f˜C
n = 4 +
−1n − 4 if n = 1 2    . 4 n n2 n a

1 −
−1n  cos
a for n = 0 1 2    , if a is not an integer. a2 − n2
 f

x sin
nxdx and integrate by parts. Write f˜S

n =

if n = m . if n = m

0

Section 15.7 1. D u
0 = D u
2 = D u
4 =

5  j=0 5  j=0 5 

cos
j ≈ −23582 D u
1 =

5 

cos
je−ij/3 ≈ 29369 − 42794i

j=0

cos
je−2ij/3 ≈ 13292 − 16579 × 10−2 i D u
3 =

D u
−4 =

cos
je−ij ≈ 96238 × 10−2 

j=0

cos
je−4ij/3 ≈ 13292 + 16579 × 10−2 i D u
−1 =

j=0

D u
−2 =

5 

5  j=0 5 

cos
je

2ij/3

cos
je

4ij/3

−2

5 

cos
jeij/3 ≈ 29369 + 42794i

j=0 5 

≈ 13292 + 16579 × 10 i D u
−3 =

cos
jeij ≈ 96238 × 10−2 

j=0 −2

≈ 13292 − 16579 × 10 i

j=0 5  1 1 −ij/3 = 245 D u
1 = e ≈ 81667 − 40415i D u
2 ≈ 65 − 17321i j + 1 j + 1 j=0 j=0 D u
3 ≈ 61667 D u
4 ≈ 65 + 17321i D u
−1 ≈ 81667 + 40415i D u
−2 ≈ 65 + 17321i D u
−3 ≈ 61667 D u
−4 ≈ 65 − 17321i 5. D u
0 = 55 D u
1 ≈ −60 + 31177i D u
2 ≈ −140 + 10392i D u
3 = 15 D u
4 ≈ −140 − 10392i D u
−1 ≈ −60 − 31177i D u
−2 ≈ −140 − 10392i D u
−3 ≈ −15 D u
−4 ≈ −140 + 10392i 5 5   7. The inverse is uj 5j=0 , where u0 = 16
1 + ik ≈ −13333 + 16667i u1 = 16
1 + ik eik/3 ≈ −42703

3. D u
0 =

5 

k=0

k=0

+54904i u2 ≈ −16346 × 10−2 + 561i u3 ≈ 33333 + 5i u4 ≈ 84968 + 27233i, and u5 ≈ 15937 − 2049i 6 6 1 1 9. u0 = e−ik ≈ 10348 + 14751 × 10−2 i u1 = e−ik e2ik/7 ≈ 93331 − 29609i 7 k=0 7 k=0 u2 ≈ −94163 × 10−2 + 88785 × 10−2 i u3 ≈ −23947 × 10−2 + 62482 × 10−2 i u4 ≈ 43074 × 10−3 + 51899 × 10−2 i u5 ≈ 25788 × 10−2 + 43852 × 10−2 i u6 ≈ 51222 × 10−2 + 34325 × 10−2 i 11. u0 ≈ −1039 u1 ≈ 42051 + 29456i u2 ≈ 13143 + 31205 × 10−2 i u3 ≈ 13143 − 31205 × 10−2 i u4 ≈ 42051 − 29456i ki cos
2 − 1 1
2 1 sin
2 + . For the DFT cos
e−ki d = − 13. The Fourier coefficients are dk = 2 0 2  2 k2 − 1 2  2 k2 − 1 127 1  approximation, choose N = 128, and approximate dk by fk = cos
j/64e−ijk/64 . 128 j=0 Then 1 d0 = sin
2 ≈ 045465 f0 ≈ 046017 2 d1 ≈ −51259 × 10−2 − 2508i and f1 ≈ −45737 × 10−2 − 25075i d2 ≈ −11816 × 10−2 − 11562i and f2 ≈ −62931 × 10−3 − 11553i d3 ≈ −51767 × 10−3 − 75984 × 10−2 i and f3 ≈ 34589 × 10−4 − 75849 × 10−2 i

Answers and Solutions to Selected Problems

A35

d−1 ≈ −51259 × 10−2 + 2508i and f−1 ≈ −45737 × 10−2 + 25075i d−2 ≈ −11816 × 10−2 + 11562i f−2 ≈ −62931 × 10−3 + 11553i d−3 ≈ −51767 × 10−3 + 75984 × 10−2 i f−3 ≈ 34589 × 10−4 + 75849 × 10−2 i
1 1 1 if k = 0, and d0 = 13 . The DFT approximation is  2 e−2ki d = +i 15. Now, dk = 2 2 2 k 2k 0  127  j 1  fk = e−ijk/64 . 128 j=0 128 d0 = 13  f0 ≈ 32944 d1 ≈ 50661 × 10−2 + 15915i f1 ≈ 46765 × 10−2 + 15912i d2 ≈ 12665 × 10−2 + 79577 × 10−2 i f2 ≈ 87691 × 10−3 + 79514 × 10−2 i d3 ≈ 5629 × 10−3 + 53052 × 10−2 i f3 ≈ 17329 × 10−3 + 52956 × 10−2 i d−1 ≈ 50661 × 10−2 − 15915i f−1 ≈ 46765 × 10−2 − 15912i d−2 ≈ 12665 × 10−2 − 79577 × 10−2 i f−2 ≈ 87691 × 10−3 − 79514 × 10−2 i d−3 ≈ 5629 × 10−3 − 53052 × 10−2 i f−3 ≈ 17329 × 10−3 − 52956 × 10−2 i

Section 15.8 1. The complex Fourier series is 2 +



  i ikt e , so S10 18 = 10207 + 16653 × 10−16 i. Approximate k k=− k=0

  127 1  1 ≈ V eik/8 , which is 10552 + 10−14 i 8 128 k=−0 k 10207 + 16653 × 10−16 i −
10552 − 20983 × 10−16 i = 0003452

 1 sin
2 i cos
2 − 1 ikt − + k e . Then 3. The complex Fourier series is 2  2 k2 − 1 2  2 k2 − 1 k=− S10

k=±1

    127 1  1 1 = 10672 − 34694 × 10−18 i. Approximate S10 ≈ V eik/8 , which works out to S10 8 8 128 k=0 k 10428 + 38025 × 10−15 i 10672 − 34694 × 10−18 i −
10428 + 38025 × 10−15 i = 0002440    1 3 1 i 2k2  2 − 3 2kit e 5. The complex Fourier series is + + , and 4 k=− k=0 4  2 k2 4  3 k3   127   1  1 ≈ V eki/2 , which works out to S10 41 = −72901 × 10−4 + 10408 × 10−17 i. Approximate S10 4 128 k=0 k 34826 × 10−3 + 91593 × 10−16 i  − 72901 × 10−4 + 10408 × 10−17 i −
34826 × 10−3 + 91593 × 10−16 i ≈ 0004212 7. 014386 − 012455i 9. −65056 × 10−3 − 2191 × 10−3 i. In the graphs associated with Problems 11 and 13, the series 1 points are the actual values computed from the transform, and the series two points are the DFT approximations. In all cases the approximations can be improved by choosing N larger. 2 sin
2 + cos
2 −  sin
 − cos
 2 cos
2 − sin
2 −  cos
 + sin
 11. fˆ
 = +i . Note: in the sum 2 2 of equation (15.24), 11 ≤ j ≤ 20 because f
t is zero outside the interval [1,2). Generate the following table using L = 4: k 1 2 3 4 5 6 7 8 9 10 11 12

fˆ
k/4 13845 − 05673i 10579 − 10421i 05761 − 13488i 02068 − 14404i −05162 − 13098i −09483 − 09865i −12108 − 05325i −12708 − 00295i −11323 + 04386i −08329 + 07966i −04359 + 09953i −00166 + 10168i

DFT fˆ
k/4 13764 − 05714i 10445 − 1048i 05558 − 13521i −0573 − 14372i −05453 − 12959i −09749 − 09602i −12288 − 04945i −12751 + 0170i −11194 + 04866i −08026 + 08393i −03909 + 10258i 0374 + 10299i

  ˆ  f
k/4 14962 1485 14667 14406 14078 13684 13227 12711 12143 11525 10866 10169

    DFT fˆ
k/4 14903 14796 14619 14383 1406 13684 13247 12752 12206 11613 10978 10306

A36

Answers and Solutions to Selected Problems The plots below compare the (a) real parts, (b) imaginary parts, and (c) absolute values, of fˆ
k/4 and the DFT approximations.

2.0 1.5

1.5

1.0

1.0 f 0.5

f

0

0


0.5

DFT


0.5


1.0


1.0
1.5

DFT

0.5


1.5 1

2

3

4

5

6 7 (a)

8


2.0

9 10 11 12

1

2

3

4

5

6 7 (b)

8

9 10 11 12

1.6 f

1.4 1.2

DFT

1.0 0.8 0.6 0.4 0.2 0

1

2

3

4

5

6

7

8

9 10 11 12

(c) cos
 − 1 sin
 +i . With L = 4 we obtain the following table (with j summing from 0 to 10 in 13. fˆ
 =   equation (15.24)):

k 1 2 3 4 5 6 7 8 9 10 11 12

fˆ
k/4 9896 − 1244i 9589 − 2448i 9089 − 3577i 8415 − 4597i 7592 − 5477i 6650 − 6195i 5623 − 6733i 4546 − 7081i 3458 − 7236i 2394 − 7205i 1388 − 6997i 0470 − 6633i

DFT fˆ
k/4 10686 − 1318i 1035 − 25925i 98047 − 37821i 90716 − 48489i 81791 − 57604i 71616 − 64909i 60574 − 70222i 49076 − 73447i 37537 − 74574i 26362 − 73678i 15924 − 70914i 0655 − 66506i

  ˆ  f
k/4 99739 98965 97675 95888 93614 90885 87722 84147 80198 75923 71333 66496

    DFT fˆ
k/4 10767 1067 10509 10286 10004 96654 92738 88334 83488 78252 7268 66828

Answers and Solutions to Selected Problems

A37

The plots below compare the (a) real parts, (b) imaginary parts, and (c) absolute values, of fˆ
k/4 and the DFT approximations.

1.2 1.0 DFT

0
0.1
0.2
0.3
0.4
0.5
0.6
0.7
0.8

0.8 f 0.6 0.4 0.2 0

1

2

3

4

5

6

7

8

9

10 11 12

1

2

3

4

5

6 7 (b)

(a)

1.2

8

9 10 11 12

DFT

1.0 f

0.8 0.6 0.4 0.2 0

1

2

3

4

5

6

7

8

9 10 11 12

(c)    i 4
2n − 1t
−1 +
−1n enit/2 , or sin . 2 n=− n=0 n n=1
2n − 1 N  nt  2
1 −
−1n  sin . The N th Cesàro sum is Then SN
t = 2 n=1 n  n  i  n  2 N  N   nt      
−1 +
−1n enit/2 , or N
t =
1 −
−1n  sin . 1−  1−  N
t = N n N n 2 n=−N n=1 Graphs are drawn for N = 5 10 25.

15. The Fourier series of f is



1.0

1.0 0.5


2

0
1 Ces`aro

Fourier

0.5


2 1

2

t


1

1

0


0.5


0.5


1.0


1.0

N=5

N = 10

2

t

A38

Answers and Solutions to Selected Problems

1.0 0.5


2
1

0

1

t

2


0.5
1.0 N = 25 N  i 2

−1n − cos
n/2enit =
cos
n/2 −
−1n  sin
nt N
t = n=−Nn=0 n n=1 n  n  i  n  2 N N        

−1n − cos
n/2enit =
cos
n/2 −
−1n  sin
nt 1−  1−  N n N n n=−Nn=0 n=1 Graphs are drawn for N = 5 10 25.

17. SN
t =

N 

1.0

1.0

0.5 –1.0


1.0

–0.5 0

0.5

1.0

t

0.5


0.5

0

–0.5


0.5

–1.0


1.0

N=5

N = 10

1.0
1.0


0.5

0.5 0

0.5

1.0

t


0.5
1.0 N = 25    N  1 1 1 −
−1n 17 i 7 n n + 1 −
−1

−1 + + − 1 enit 4 n=−Nn=0 2 n2  2 n 2 2   N 5 17  2 1 −
−1n n + − + 3
−1 sin
nt = cos
nt − 4 n=1 n2  2 n 2    n  1 1 −
−1n N   1 17 i 7   n n N
t = + 1 −
−1


−1 1−  + + − 1 enit 4 n=−Nn=0 N 2 n2  2 n 2 2  

 n  1 −
−1n N  5 17  2   n + − + 3
−1 = cos
nt − sin
nt 1−  4 n=1 N n2  2 n 2 Graphs are drawn for N = 5 10 25 (see next page).

19. SN
t =

0.5

1.0

t

Answers and Solutions to Selected Problems

7

7

6

6 5

5

4

4 3

3

2

2

1
1.0


0.5

A39

0 N=5

0.5

1

t

1.0


1.0


0.5

0 N = 10

0.5

1.0

t

7 6 5 4 3 2 1
1.0


0.5

0 N = 25

0.5

1.0

t

N 

N  i −2
3
−1n − 1enit/2 = 1 +
3
−1n − 1 sin
nt/2 N
t = n n n=−Nn=0 n=1 N     −2 1 −  Nn 
3
−1n − 1 sin
nt/2 1+ n

21. SN
t = 1 +

n=1

Hamming filtered partial sum: HN
t = 1 + N 

N  −2
054 + 046 cos
n/N
3
−1n − 1 sin
nt/2 n n=1

2 −a 2 n2 /N 2 e
3
−1n − 1 sin
nt/2 n n=1 1 Graphs are drawn for N = 5 10 25, with a = in the Gauss filter. 2

Gauss filtered partial sum: GN
t = 1 +

−

4

Fourier

3

2 Ces`aro
2


1

Gauss

Ces`aro

1 0

1

2

t


2


1


2 Hamming

Gauss

2

1
1

Fourier

4

3

N=5


1

0


2 Hamming

N = 10

1

2

t

Answers and Solutions to Selected Problems

A40

Fourier 4 3

Gauss

2 Ces`aro
2


1

1
1

0


2 Hamming

N = 25

Section 15.9 1. Power spectrum of y
t = 4 sin
80t − sin
20t

10 4 10 3 10 2 10 1 10 0 10
1 10
2 10
3

0

50

100 150 200 250 300 350 400 450 500

3. Power spectrum of y
t = 3 cos
90t − sin
30t

10 4 10 3 10 2 10 1 10 0 10
1 10
2 10
3

0

50

100 150 200 250 300 350 400 450 500

1

2

t

Answers and Solutions to Selected Problems 5. Corrupted signal of y
t = cos
30t + cos
70t + cos
140t

6 4 2 0
2
4
6
8

0

5

10

15

20

25

30

35

40

45

50

Frequency components of the corrupted signal

160 140 120 100 80 60 40 20 0

50 100 150 200 250 300 350 400 450 500

7. Corrupted signal of y
t = cos
20t + sin
140t + cos
240t

6 4 2 0
2
4
6
8

0

5

10

15

20

25

30

35

40

45

50

A41

Answers and Solutions to Selected Problems

A42

Frequency components of the corrupted signal

150

100

50

0

50 100 150 200 250 300 350 400 450 500

CHAPTER 16 Section 16.1 3. We will illustrate for P3
x. By Rodrigues’s formula, 1 d3 1 1

x2 − 13  =
120x3 − 72x =
5x3 − 3x. P3
x = 3 2 3! dx3 48 2 1 
6 − 2k! 3 = 1 and x3−2k =
−1k 5. For example, with n = 3 we have 2 8k!
3 − k!
3 − 2k! k=0 5 4! 3 6! x3 − x = x3 − x = P3
x 8
3!
3! 8
2!
1! 2 2 After cancellations because Pn
x is one solution, we obtain 7. Substitute 
Q
x = Pn
xu
x in Legendre’s equation. P

x Pn v

x 2x 2x



u + 2 − = −2 n + u = 0. Let v = u to write . Integrate to obtain 2 Pn 1 − x v
x Pn
x 1 − x2
1 1 . Then u
x = dx and ln
v
x = −2 ln
Pn
x − ln
1 − x2 , so v
x = Pn
x2
1 − x2  Pn
x2
1 − x2  Q
x = Pn
xu
x. −112 2 +  4 + 1008 1 12  2 − 10 1
5x3 −3x+660
63x5 −70x3 +15x+· · · . 13. For −1 < x < 1 sin
x/2 = 2 x +168 4   2 6 8 In the graph, the sum of these terms is indistinguishable from the function.

y 1.0 0.5
1.0


0.5

0
0.5
1.0

0.5

1.0

x

Answers and Solutions to Selected Problems

A43



1 5 15 15 1 1 + − cos
1 sin
1 + − cos2
1
3x2 − 1 + 15. For −1 < x < 1 sin2
x = − cos
1 sin
1 + 2 2 8 8 4 2

531 585 585 1 cos
1 sin
1 − + cos2
1
35x4 − 30x2 + 3 + · · · . In the graph, the sum of these terms is 32 32 16 8 indistinguishable from the function.

y 0.7 0.6 0.5 0.4 0.3 0.2 0.1
1.0


0.5

0

0.5

1.0

x

71 3 11 1
5x3 − 3x +
63x5 − 70x3 + 15x + · · · . The graph below 17. For −1 < x < 0 and for 0 < x < 1 f
x = x − 2 82 16 8 shows the function and the sum of these three terms of the Fourier–Legendre expansion.

y 1.0 0.5
1.0


0.5

0 0.5
0.5

1.0

x


1.0

Note that in 13 and 15, just the first three terms of the eigenfunction expansion approximate the function so closely that the two graphs are virtually indistinguishable. In 17, it is clear that we must include more terms of the eigenfunction expansion to reasonably approximate f
x.

Section 16.2 1. Let y = xa J
bxc  and compute y
= axa−1 J
bxc  + xa bcxc−1 J

bxc  and Y

= a
a − 1xa−2 J
bxc  + 2axa−1 bcxc−1 + xa bc
c − 1xc−2 J

bxc  + xa b2 c2 x2c−2 J


bxc . Substitute these into the differential equation and simplify to get c2 xa−2 
bxc 2 J


bxc  + bxc J

bxc  +
bxc 2 −  2 J
bxc  = 0. 3. y = c1 J1/3
x2  + c2 J−1/3
x2  5. y = c1 x−1 J3/4
2x2  + c2 x−1 J−3/4
2x2  7. y = c1 x4 J3/4
2x3  + c2 x4 J−3/4
2x3  9. y = c1 x−2 J1/2
3x3  + c2 x−2 J−1/2
3x3  √ √ 11. y1 = c1 J3
x + c2 Y3
x 13. y = c1 J4
2x1/3  + c2 Y4
2x1/3  15. y = c1 x2/3 J1/2
x + c2 x2/3 J−1/2
x

 


+ zJ1/3 + z2 − 19 J1/3 = 0 17. Substitute into the differential equation and use the face that J1/3
z satisfies z2 J1/3 √ √ 19. y = c1 xJ2
x + c2 xY2
x 21. y = c1 x2 J2
x + c2 x2 Y2
x 29. (a) The sum of the first five terms of the Fourier–Bessel expansion is ≈ 167411J2
5135x − 077750J2
8417x + 08281J2
11620x − 06201J2
14796x + 06281J2
17960x. From the graphs of x and the first five terms of this expansion, more terms would have to be included to achieve reasonable accuracy.

Answers and Solutions to Selected Problems

A44

(c) The sum of the first five terms of the Fourier–Bessel expansion is 085529J2
5135x − 021338J2
8417x +035122J2
11620x − 020338J2
14796x + 0025800J2
17960x. As the graphs indicate, more terms are needed to reasonably approximate the function by the partial sum of the Fourier–Bessel expansion.

y

y

1.0

0.4

0.8

0.3

0.6 0.2 0.4 0.1

0.2 0.2

0

0.4

0.6 (a)

0.8

1.0

x

x 0

0.2

0.4

0.6

0.8

1.0

(c)

Section 16.3

2  2n − 1  for n = 1 2    ; eigenfunctions nonzero constant multiples of 1. regular problem on [0, L]; eigenvalues 2L   2n − 1 sin x 2L     2n − 1  2n − 1  2 x  cos 3. regular on [0, 4]; 2 4 2 4 1 5. periodic on [−3 3]; 0 is an eigenvalue with eigenfunction 1; for n = 1 2     n2 is an eigenvalue with 9 eigenfunction an cos
nx/3 + bn sin
nx/3, not both an and bn zero √ 1√ . There are infinitely many such solutions, of 7. regular on [0, 1]; eigenvalues are positive solutions of tan
 = 2 which the first four are approximately 0.43, and 89.82. Corresponding to an eigenvalue n , √ √  10.84, 40.47, eigenfunctions have the form 2  cos
n x + sin
x. 9. regular on 0  1 + n2 for n = 1 2     e−x sin
nx for n = 1 2     n n2  2 11. regular on 1 e3  1 + for n = 1 2     x−1 sin ln
x for n = 1 2    9 3  2 sin
nx. The tenth partial sum of the series is compared to the function in the graph. 13. For 0 < x < 1 1 − x = n=1 n

y 1.0 0.8 0.6 0.4 0.2 0

0.2

0.4

0.6

0.8

1.0

x

Answers and Solutions to Selected Problems

A45

√ √    4 2 cos
n/2 − 2 sin
n/2 −
−1n 2n − 1 cos x . This converges to −1 for 0 < x < 2 15. The expansion is 2n − 1 8 n=1  and 1 for 2 < x < 4, and to 0 if x = 2. The tenth partial sum of the series is compared to the function in the graph.

y 1.0 0.5 0

1

2

3

4

x


0.5
1.0 17. For −3 < x < 3 x2 = 3 2 + 36


−1n cos
nx/3 The tenth partial sum of the series is compared to the n2 n=1

function in the graph.

y 80 60 40 20
8
6
4
2

2 4

0

6

8

x

   
4 1 1 2n − 1 2n − 1 x  Now,
f · n  = x dx = 19. Normalized eigenfunctions are n
x = √ cos x
4 − x √ cos 8 8 0 2  2  √ 4
−1n +
2n − 1 √ 4
−1n +
2n − 1 2 4 2  512
128 2 −128 2 , so ≤ f · f = 0 x
4 − x2 dx = 15 ,  3
2n − 13  3
2n − 13 n=1   4
−1n +
2n − 1 2 512 1  . ≤ = or √  3
2n − 13 960 n=1 15
128 22

Section 16.4 3. The interval on which 13
t is nonzero is disjoint from the interval on which −21
t is nonzero, so 13
t−21
t is identically zero, hence − 13
t−21
tdt = 0. Graphs of 13
t and −21
t are shown below.

y 1.0 0.5
10


5

0
0.5
1.0

5

10

t

A46

Answers and Solutions to Selected Problems

5. A graph of 
2t − 3 is shown below.

y 1.0 0.5 2
3
2
1 0

1

3

4

t

5


0.5
1.0 7. The Fourier 16 is  seriesof f
t  on −16      n   6 3n n 3n 4 −2 sin + sin − sin − 2 sin + 4 2 n 8 16 n=1 n         n   nt 4   3n n 6 n 3n sin cos + − − cos
n + 2 cos − cos + cos − 2 cos + 4   16 n 4 2 n 8 16  n nt sin . The graph below compares the function with the fiftieth partial sum of its Fourier series. cos 4 16

y
6
4
2
15
10
5

5


2

10

15

t


4
6    n  n  nt  1 3n −14 sin − 8 sin + 16 sin cos + 9. The Fourier series of f
t on [−16, 16] is 2 4 8 16 n=1 n      1 n 3n nt n 3
−1n + 2 cos + 3 + 8 cos − 16 cos sin . The graph below compares the function n 2 4 8 16 with the fiftieth partial sum of its Fourier series.

y 8 6 4 2
15
10
5
2
4
6
8

5

10

15

t

Answers and Solutions to Selected Problems

A47

CHAPTER 17 Section 17.1

2 y

2 y2 n2  2 c2 n2  2 = − = − sin
nx/L cos
nct/L and sin
nx/L cos
nct/L.

x2 L2

t L2 2 2

y 1

y 1 3. Compute 2 = f


x + ct + f


x − ct and 2 = c2 f


x + ct + f


x − ct.

x 2

t 2  2 

2 z

z 2 z 5. The problem for the displacement function z
x y t is 2 = c2 + for 0 < x < a 0 < y < b

t

x2 y2 1. Compute

z
x y 0 = f
x y

z
x y 0 = 0 for 0 < x < a 0 < y < b

t

z
0 y t = z
a y t = z
x 0 t = z
x b t = 0

Section 17.2 1. 3. 5. 7.

     2
−1n 16
−1n
2n − 1x
2n − 1ct sin sin + sin
nx sin
nct y
x t = 3 3 2 2 2 2 n=1
2n − 1  c n=1 n  c      2
2n − 1t 108
2n − 1x sin y
x t = sin 44
2n − 1 3 3 n=1    √  24
2n − 1x
−1n+1 sin cos
2n − 1 2t y
x t = 2 2 n=1
2n − 1       −32 3
2n − 1t
2n − 1x cos sin y
x t = 3 3 2 2 n=1
2n − 1          4 n n 3nt nx + cos − cos sin sin 2 2 2 4 2 2 n=1 n  

9. Let Y
x t = y
x t + h
x and substitute into the problem to choose h
x = 19 x3 − 49 x. The problem for Y becomes

2 Y

2 Y = 3 

t2

x2 Y
0 t = Y
2 t = 0 4 Y 1
x 0 = 0 Y
x 0 = x3 − x 9 9 t  √   nx  32 n 3t n cos , and then y
x t = Y
x t − h
x.
−1 sin We find that Y
x t = 33 3n 2 2 n=1 11. Let Y
x t = y
x t + h
x and choose h
x = cos
x − 1. The problem for Y is

2 Y

2 Y = 2 2

t

x Y
0 t = Y
2 t = 0

Y
x 0 = 0

t      16 1
2n − 1x
2n − 1t sin cos , and then This problem has solution Y
x t = 2 2 2 n=1 
2n − 1
2n − 1 − 4 y
x t = Y
x t + 1 − cos
x.  r t  nx 1 r t  nx  2A
L 13. u
x t = e−At/2 rn cos n + sin n , where Cn = dx and Cn sin f
x sin L AL 2L 2L rn 0 L n=1  rn = 4
BL2 + n2  2 c2  − A2 L2 Y
x 0 = cos
x − 1

Answers and Solutions to Selected Problems

A48

    64 1 nx 2 1 sin 15. (a) The solution with the forcing term is yf
x t = − 3 3 2
2n − 1 9
2n − 1
2n − 1 − 16 4 n=1    3
2n − 1t 1 cos +

cos
x − 1. 4 9 2      128 3
2n − 1t
2n − 1x cos . sin (b) Without the forcing term, the solution is y
x t = 3 3 4 4 n=1 
2n − 1 Both solutions are graphed together for times t = 05 06 49, and 9.8, using the same axes to allow comparison.

t = 4.9 2

t = 0.5

1

t = 0.6

0 0

1

2 x

3

4

–1

t = 9.8 –2

17. yj1 = 025 for j = 1 2     19 y12 = 008438 yj2 = 005 for j = 2 3     19 y13 = 013634 y23 = 0077149 yj3 = 0075 for j = 3 4     19 y14 = 017608 y24 = 010786 y34 = 010013 yj4 = 01 for j = 4 5     19 y15 = 020055 y25 = 014235 y35 = 012574 y45 = 012501 yj5 = 0125 for j = 5 6     19 19. We give yjk for j = 1 2     9, first for k = −1, then k = 0 1     5. yj−1  008075 0127 014475 014 011875 0087 005075 0016 −001125 yj0  0081 0128 0147 0144 0125 0096 0063 0032 0009 yj1  0079125 01735 014788 0147 013063 010475 0075375 00485 0030125 yj2  00057813 002115 077078 014903 013567 011328 0087906 0065531 0050516 yj3  −0055066 027160 13199 018908 014015 012162 010062 0083022 0068688 yj4  −0092055 03768 17328 029675 014653 012981 011355 010072 0083463 yj5  −0093987 053745 19712 048652 016125 013803 012669 011814 00941

Section 17.3


−1 sin
 10 cos
x cos
12td 3. y
x t = sin
x sin
4td 
25 + 2  2 2 − 1 0 0   


 1 2 cos
 −  sin
 1 −2  cos
 + 2 sin
 cos
x + sin
x sin
3td 5. y
x t = e−2 e 3 4 + 2 3 4 + 2 0
2 2 −  sin
 − 2 cos
 sin
x cos
3td 7. y
x t =  3 0 1. y
x t =






Answers and Solutions to Selected Problems


1 sin
/2 − sin
5/2 sin
x sin
2td  2 − 1
3 −

16 cos3
 − 12 cos
 + 122 sin
 cos2
 − 32 sin
 − 8 sin
 cos2
 11. y
x t = 75 0 + 2 sin
 + 2 sin
x sin
14td  x x  13. y
x t = At +
1 − A t − H t− c c 9. y
x t =



0

Section 17.4 1. Characteristics are lines x − t = k1  x + t = k2 
x+t

 − d = x2 + t2 − xt y
x t = 21
x − t2 +
x + t2 + 21 x−t

3. Characteristics are x − 7t = k1  x + 7t = k2  49 1 y
x t = cos

x − 7t + cos

x + 7t + t − x2 t − t3 2 3 5. Characteristics are x − 14t = k1  x + 14t = k2  y
x t = 21 ex−14t + ex+14t  + xt 7. y
x t = x + 18
e−x+4t − e−x−4t  + 21 xt2 + 16 t3 9. y
x t = x2 + 64t2 − x + 321
sin
2
x + 8t − sin
2
x − 8t + 121 t4 x 11. y
x t = 21 cosh
x − 3t + cosh
x + 3t + t + 41 xt4 In each of 13, 15, and 17, the graphs show a progression of the motion as a sum of forward and backward waves. 13.

u

u

1.0 0.6 0.4

0.5

0.2
4


2

0

2

4

x


4


0.5


2

0
0.2

2

x

4


0.4
0.6


1.0 (2)

(1)

u

u 0.6

0.2

0.4 0.4

0.2
4


2

0
0.2
0.4
0.6 (3)

2

4

x


4


2

0
0.2
0.4 (4)

2

4

x

A49

Answers and Solutions to Selected Problems

A50

u 0.4 0.2
4


2

0

2
0.2

x

4


0.4 (5) 15. u

u

1.0 0.8

0.8

0.6

0.6


4


2

0.4

0.4

0.2

0.2

0 (1)

2

4

x


4


2

0 (2)

x

2

4

2

4

u

u 0.7

0.5

0.6

0.4

0.5 0.4

0.3

0.3

0.2

0.2 0.1

0.1
4


2

0 (3)

2

4

x


4


2

0 (4)

u

u


4


2

0.5

0.5

0.4

0.4

0.3

0.3

0.2

0.2

0.1

0.1

0 (5)

x

2

4

x


8
6
4
2

0 (6)

2

4

6

8

x

Answers and Solutions to Selected Problems 17.

u

u
4


2


4


2

0

2

0

x

4


4


2

2


0.5


0.4


1.0


0.8


1.5


1.2


2.0


1.6

(1)

(2)

u

u

0
0.2

2

x

4


4


2

0
0.2


0.4


0.4


0.6


0.6


0.8


0.8


1.0


1.0

(3)

2

4

4

x

x

(4)

Section 17.5 2

1

xJ0
2405xdx 01057 = 078442, a2 = 004112, a3 = −81366, =2

J1
24052 02695 a4 = −3752, a5 = −64709. The fifth partial sum of the series gives the approximation z
r t ≈ 078442J0
2405r cos
2405t + 004112J0
5520r cos
5520t − 81366J0
8654r cos
8654t− 3752J0
11792r cos
11792t − 64709J0
14931r cos
14931t. The graph below shows z
r t at various times.

1. We find that (approximately) a1 =

0

z 4000 2000

t  2.7 t  1.4 t  0.3 0.2


2000
4000

0.4

0.6

0.8

1.0 r

0

t  0.9

A51

A52

Answers and Solutions to Selected Problems

3. Approximately, z
r t ≈ 12534J0
2405r cos
2405r − 088824J0
5520r cos
5520t− 2489J0
8654r cos
8654t − 11336J0
11792r cos
11792t − 19523J0
14931r cos
14931t. The graph below shows z
r t at selected times.

z 15,000 10,000 5,000 0
5,000
10,000
15,000
20,000

t  1.5 t  2.7 0.4 0.2

0.8 0.6

1.0

r

t  0.9 t  0.4 t  4.2

Section 17.6 1. Compute 1
 1 1 0
r =
4 − r 2  sin2
d =
4 − r 2  2 2 − 2 ⎧ 0 ⎪ ⎨ 2
 4−r   sin2
 cos
nd = 4 − r 2 n
r = 1 ⎪  − −  ⎩  2

if n = 2 if n = 2

4 − r2
 2 sin
 sin
nd = 0  − 
1     2 1 2 Thus, z
r  t = j 
1 −  J
j d J r cos
j0k t 0 0k 0 0k

J1
j0k 2 0 2 k=1 
1     4 1 2 + j 
 − 1J
j d J r cos
2 cos
j2k t 2 2k 2

J3
j2k 2 0 2 2k k=1  
1   4
−1p+1 1 + j sin
p J
j dJ r sin
jpq t p pq  pq 2 2 0 p=1 q=1 pjpq Jp+1
jpq  n
r =

≈ 11081J0
12025r cos
240483t − 013975J0
2760r cos
552008t +04555J0
43270r cos
865373t − 002105J0
58960r cos
117915t +001165J0
74655r cos
1443092t + · · · − 29777J2
25675r cos
2 cos
51356t −14035J2
42085r cos
2 cos
841724t − 11405J2
58100r cos
2 cos
116198t −083271J2
7398r cos
2 cos
147960t − · · · 

Section 17.7

 

  1 8
−1n+1  2 16 1  nx n 2 + 3
−1 − 1 sin sin
y cos 1. z
x y t = n + 4t  n=1 n n 2 2

Answers and Solutions to Selected Problems      16
2n − 1x
2m − 1x cos 3. z
x y t = sin  2 2 2 2 2 n=1 m=1 
2n − 1
2m − 1
2n − 1 +
2m − 1 
2n − 12 +
2m − 12 t × sin  



CHAPTER 18 Section 18.1

2 u

u

u = k 2 for 0 < x < L t > 0; u
0 t =
L t = 0 for t ≥ 0 u
x 0 = f
x for 0 ≤ x ≤ L.

t

x

x

u

2 u

u 3. = k 2 for 0 < x < L t > 0
0 t = 0 and u
L t = 
t for t ≥ 0; u
x 0 = f
x for 0 ≤ x ≤ L.

t

x

x 1.

Section 18.2 In these solutions, exp
A = eA .

    8L2 −
2n − 12  2 kt
2n − 1x exp sin 3 3 L L2 n=1
2n − 1       −16L
2n − 1x −3
2n − 12  2 t sin exp u
x t = 2 L L2 n=1
2n − 1
2n − 1 − 4  2 4 2 cos
nxe−4n t u
x t =  2 − 2 3 n n=1    nx 2 2  1 1 − e−6
−1n e−n  t/18 cos u
x t =
1 − e−6  + 12 2 2 6 36 + n  6 n=1      4B
2n − 1x −
2n − 12  2 kt sin exp u
x t = 2L 4L2 n=1
2n − 1

1. u
x t = 3. 5. 7. 9.



11. Substitute ex+t v
x t into the partial differential equation and solve for  and  so that v1 = kvxx . We get  = −A/2 and  = k
B − A2 /4. 13. Let u
x t = e−3x−9t v
x t. Then vt = vxx  v
0 t = v
4 t = 0, and v
x 0 = e3x . Then v
x t =   2 2   1 − e12
−1n nx e−n  t/16 2n sin . Graphs of the solution are shown for times t = 0003 002 008, 144 + n2  2 4 n=1 and 13.

u

u

400

4

200 0
200

2 1

2

3

4

x

0
2


400


4


600


6


800

1

2


8 t = 0.003

t = 0.02

3

4

x

A53

Answers and Solutions to Selected Problems

A54

u

u

1.0

0.0016 0.0014 0.0012 0.0010 0.0008 0.0006 0.0004 0.0002

0.8 0.6 0.4 0.2 1

0

2 t = 0.08

3

4

x

0

1

2 t = 1.3

3

x

4

15. Let u
x t = v
x t + f
x and choose f
x = 3x + 2 to have vt = 16vxx , v
0 t = v
1 t = 0, and   2 2  4n 
−1n + 2
−1n − 2 − 2n2  2 2 2 sin
nxe−16n  t and u
x 0 = x2 − f
x. Then v
x t = 2 n3  3 n=1 u
x t = v
x t + 3x + 2. Graphs of the solution for times t = 0005 0009, and 0.01.

u

u

5.0

5.0

4.5

4.5

4.0

4.0

3.5

3.5

3.0 2.5

3.0

2.0

2.5

1.5 0

0.2

0.4 0.6 t = 0.005

0.8

1.0

x

2.0 0.2

0

0.4 0.6 t = 0.009

u 5.0 4.5 4.0 3.5 3.0 2.5 2.0 0

0.2

0.4

0.6 t = 0.01

0.8

1.0

x

0.8

1.0

x

Answers and Solutions to Selected Problems

A55

17. Let u
x t = e−At w
x t. Then wt = 4wxx  w
0 t = w
9 t = 0, and w
x 0 = 3x. Obtain  nx  54
−1n+1 2 2 w
x t = sin e−4n  t/81 . The graphs compare solutions at times t = 0008 004, and 0.6 for n 9 n=1 A = 41 , 1, and 3.

u

u

A = 14

25

20

A=1

20

A=1

15

A=3

15

A = 14

A=3

10

10

5

5 x

0

2

4 6 t = 0.008

x

8

2

0

4 6 t = 0.04

8

u A = 14 10 8 6

A=1

4 A=3

2

x 2

0

4 6 t = 0.6

8

  1 1 −
−1n  1 −
−1n 4  2 2 −4n2 t sin
nxe−4n t . The graphs compare −1 + 4n sin
nx + t + e n5  n=1 n3 n=1 8 solutions with and without the source term, for times t = 08, 04, and 13.

19. u
x t =





u

u

0.25

0.5

With source

0.20

0.4

0.15

0.3

No source 0.10

0.2

0.05

0.1

0

0.5

1.0

1.5 2.0 t = 0.8

2.5

3.0

x

0

0.5

1.0

1.5 2.0 t = 0.4

2.5

3.0

x

Answers and Solutions to Selected Problems

A56

u 0.30 0.25 0.20 0.15 0.10 0.05 0

0.5

1.0

1.5 2.0 t = 1.3

2.5

x

3.0

 nx  50 1 −
cos 5
−1n  2 2 −n2  2 t/25 −25 + n sin 21. u
x t =  t + 25e n3  3 n2  2 − 25 5  n=1   nx 2 2  4
−1n + 2 e−n  t/25 . The graphs compare solutions with and without the source term, for + −250 sin n3  3 5 n=1 times t = 07, 15, 26, and 42.

u

u

No source

No source

12

8

10

Source

Source

6

8 6

4

4 2

2 1

0

2

3 t = 0.7

4

5

x

0

1

2

3 t = 1.5

4

x

5

u

u No source

3

5

No source

2

4 3

Source

Source

1

2 0

1 0

1

2

3 t = 2.6

4

5

x

1

2

3


1 t = 4.2

4

5

x

Answers and Solutions to Selected Problems  27
−1n+1 23. u
x t = 128 n=1



2  2 t/9

16n2  2 t + 9e−16n n5  5

−9

 sin

 nx 3

+ 2K

    1 −
−1n nx −16n2  2 t/9 sin e . n 3 n=1

The graphs show the solution with and without the source term, at times t = 005 and 02, with K = 21 .

u

u 0.035

0.25

0.030

0.20

0.025

0.15

0.020

Source No source

0.015

0.10

0.010 0.05

0.005 x

0

0.5

1.0

1.5 2.0 t = 0.05

2.5

3.0

x 0

0.5

1.0

1.5 2.0 t = 0.2

2.5

3.0

25. In the following, j = 1 2     9 uj0 : 0.009, 0.032, 0.063, 0.096, 0.125, 0.144 0.147. 0.128, 0.081 uj1 : 0.0125, 0.034, 0.0635, 0.095, 0.1225 0.14, 0.1415, 0.121, 0.0725 uj2 : 0.01475, 0.064125, 0.089, 0.094, 0.1195 0.136, 0.136, 0.114, 0.0665 uj3 : 0.023381, 0.058, 0.084031, 0.099125, 0.11725 0.13188, 0.1305, 0.10763, 0.06175 27. In the following, j = 1 2     9. uj0 : 0.098769, 0.19021, 0.2673, 0.32361, 0.35355 0.35267, 0.31779, 0.24721, 0.14079 uj1 : 0.096937, 0.18622, 0.26211, 0.31702, 0.34585 0.34417, 0.30887, 0.23825, 0.13220 uj2 : 0.095124, 0.18307, 0.25697, 0.3105, 0.33822 0.33577, 0.30004, 0.22939, 0.12566 uj3 : 0.09330, 0.17956, 0.25188, 0.30405, 0.33062 0.32745, 0.29131, 0.22112, 0.12018

Section 18.3

1
8 2 cos
xe− kt d  0
16 + 2  
 8 cos3
 − cos
 + 4 sin
 cos2
 − 2 sin
 cos
 cos
x u
x t = 2 0 

  4 −2 sin
 cos3
 + sin
 cos
 + 8 cos4
 − 8 cos2
 +  2 sin
x e− kt d −  2 2
 2 sin
xe− kt d u
x t =  0 2 + 2 2
1 − cos
h 2 sin
xe− kt d u
x t =  0   4
2 2 sin
xe− t e−t /2 d u
x t =  0
1 + 2 2  
t x 2
t − erfc √ u
x t = d, in which erfc is the complementary error function. 0 2 k

1. u
x t = 3.

5. 7. 9. 11.

A57

Answers and Solutions to Selected Problems

A58

Section 18.4

 
1 2 2 2  J
j d J0
jn re−jn t ; the fifth partial sum, with approximate values inserted, is 0 n 2

J
j  0 1 n n=1 U
r t ≈ 8170J0
2405re−5785t − 11394J0
5520re−3047t + 07983J0
8654re−7489t − 0747J0
11792re−13904t + 06315J0
14931re−22293t . A graph of this function is shown for times t = 0003 0009 004, and 07.

1. U
r t =



U t  0.003

0.8 0.6

t  0.009

0.4 t  0.04

0.2 t  0.7 0.2 0.4

0

0.6

0.8

x

1.0

   
1 2 jn 2 2 r e−jn t/18 ; the fifth partial sum, with approximate values 
9 −  J
j d J 0 n 0 2

J
j  3 0 1 n n=1 inserted, is U
r t ≈ 99722J0
2405r/3e−578t/18 − 1258J0
5520r/3e−3047t/18 + 04093J0
8654r/3e−7489t/18 − 01889J0
11792r/3e−13904t/18 + 01048J0
14931r/3e−22293t/18 . A graph of this partial sum is shown for times t = 0003 0009 008, and 04.

3. U
r t =

u



t = 0.003 t = 0.009

8 6 t = 0.08 4 t = 0.4 2 r 0

0.5

1.0

1.5

Section 18.5 1. u
x y t =

  n=1 m=1
K
L

bnm sin

2.0

2.5

 nx L

sin

3.0

 my K

 e−nm kt , where nm =

n2 m2 + L2 K 2

4 f
x y sin
nx/L sin
my/Kdx dy LK 0 0  4 8
−1m+1 m 2 sin
x sin
mye−
1+m t . 3. u
x y = 2
2m − 12 
2m + 1 m=1 nm =

CHAPTER 19 Section 19.1 1. & 2
f + g =
f + gxx +
f + gyy =
fxx + fyy  +
gxx + gyy  = & 2 f + & 2 g and & 2
f =
fxx +
fyy = 
fxx + fyy  = & 2 f .

  2 and

Answers and Solutions to Selected Problems  

y 2 − x2

2x 2x ln
x2 + y2  = 2 =2 2 and . 2 2 2

x  x +y

x x + y
x + y2 2  x2 − y 2

2 2x =2 2 . Then & 2 ln
x2 + y2  = 0, provided that x2 + y2 = 0. Similarly, 2 2 2

y x + y
x + y2 2

2 u 1 u 1 2 u 5. Recall that, in polar coordinates, Laplace’s equations is 2 + + 2 2 = 0. It is routine to verify by

r r r r  substitution that the given functions are harmonic. 3. Compute

Section 19.2 1. u
x y = 3. u
x y =

−1 sin
x sin h

y −  sinh
 2   32 n
−1n+1

sin
nx sinh
ny  2 sinh
4n
2n − 12
2n + 12  nx  ny  1 1
−1n − 1 sin
x sinh
y + sin sinh 16n 5. u
x y = sinh
 2   2
n − 22
n + 2 sinh
n 2 /2 2 2 n=1n=2 1 sin
y sinh
x + sinh
2     
2n − 1y
2n − 1x sinh , 7. u
x y = cn sin 2a 2a n=1  
a 2
2n − 1x where cn = dx. f
x sin a sinh
2n − 1b/2a 0 2a    −1 1 −
−1n 2 9. u
x y = sin
y sinh

x − 4 + 2 sin
ny sinh
nx sinh
4  3 n3 n=1 sinh
4n n=1

Section 19.3

 n  1 1 r 3. u
r  =  2 + 2
−1n 2 2 cos
n + n sin
n 3 2 n n=1  n − 1 1  r e
−1n

− cos
n − n sin
n + e2 cos
n + e2 n sin
n 5. u
r  = sinh
 +   n=1 4 n2 + 1   n   2 r

n2  2
−1n sin
n − 6
−1n sin
n 7. u
r  = 1 + 3 8 n n=1

1. u
r  = 1

9. In polar coordinates, the problem is & 2 U
r  = 0 for r < 4 U
4  = 16 cos2
. This has solution U
r  = 8 + r 2
cos2  − 21 . In rectangular coordinates, the solution is u
x y = 21
x2 − y2  + 8. 11. In polar coordinates, the solution is U
r  = r 2
2 cos2  − 1, so u
x y = x2 − y2 .

Section 19.4

   3
2  3 1  = d = 98696/; u  /3 ≈ 4813941647/, 2 8 0 5/4 − cos
 −  4 u
02 /4 ≈ 8843875590/ 3. u
4  ≈ 15525/ u
12 3/2 ≈ 302/ u
8 /4 ≈ 11156/ u
7 0 ≈ 24851/ Rn 1
2 R2 − R2 /4 5. With u
r  = r n sin
n, compute u
R/2 /2 = n sin
n/2 = 2 2 2 2 0 R + R /4 − R2 cos
 − /2 Rn sin
nd. Upon dividing out common powers of R and solving for the integral, we obtain
2 sin
n 1 2 sin
n/2 = d. n 2 3 5 − 4 sin
 0
2  n
2   1 1 7. cos = cos
nd
−1n = cos
nd n−1 n−1 3
2  2 5 − 4 sin
 3
2  5 + 4 cos
 0 0 

1. u

A59

A60

Answers and Solutions to Selected Problems

Section 19.5

    1 4−x 4+x arctan − arctan for − < x <  y > 0  y y   y
1 1 u
x y = e− cos
d −  0 y2 +
 − x2 y2 +
 + x2     2

2

u
x y = f
 sin
d sin
xe−y d + g
 sin
d sin
ye−x d  0  0 0 0  h 2 u
x y = Be−y sin
x +
1 −
−1n 
1 − e−ny  sin
nx 3 n=1  n 


−1n −ny 4
−1n 2  e −2 sin
nx. − +6 Using a finite Fourier sine transform in x, we get u
x y =  n=1 n n n   2
 sin
ye−x d u
x y =  0 1 + 2     y
8 A A x−8 x−4 − arctan + arctan u
x y = d =  4 y2 +
 − x2  y y

1. u
x y = 3. 5. 7. 9. 11. 13.

Section 19.6

 4
−1n+m  √ sin
nx sin
my sinh
 n2 + m2 z n=1 m=1 nm 2 sinh  n2 + m2    16 3. u
x y z =  2
2n − 1
2m − 1 sinh
2
2m − 12 +  2
2n − 12   n=1 m=1   × sin

2n − 1x sin

2m − 1z sinh

2m − 12 +  2
2n − 12 y ⎡ 1. u
x y z =

 

 ⎢ ⎢  16 ⎢ + ⎢  2
2n − 1
2m − 1 n=1 m=1 ⎣

  sinh 

1
2m − 12 +  2
2n − 12 4



⎤  ⎥   ⎥
2m − 12
2m − 1y 2 2 × sin

2n − 1x sin sinh + 
2n − 1 z ⎥ ⎥ 2 4 ⎦

Section 19.7 1. u
  =

  
2n + 1A
1  n
arccos
2 Pn
d Pn
cos
 2 R −1 n=0

≈ 29348A − 37011A   3

 R

P1
cos
 + 11111A

  2 R

P2
cos


  4   5 P3
cos
 + 03200A P4
cos
 − 02120 P5
cos
 + · · · R R R    2 3. u
  ≈ 60784 − 98602 P
cos
 + 52360 P2
cos
 R 1 R   3   4   5 −24044 P3
cos
 + 15080 P4
cos
 − 09783 P5
cos  + · · · R R R

1 1 5. u
  = R −1
T R  R2 − R1 1 1  2 −05397A

Answers and Solutions to Selected Problems

A61

1

AP2n−1
xdx
4n − 1A
1 = P2n−1
xdx. 1 R2n−1 0 R2n−1 0
P2n−1
x2 dx n=1  4n − 1
1 9. u
  = a2n−1 2n−1 P2n−1
cos
, where a2n−1 = 2n−1 f
arccos
xP2n−1
xdx R 0 n=1 7. u
  =



2n−1

a2n−1 

P2n−1
cos
, where a2n−1 =

0

Section 19.8

3.

5. 7. 9. 11.

1

4 cos
x cosh

1 − y + C. 4 cos
xdx = 0, a solution may exist. We find u
x y = − sinh
   Since 0 cos
3xdx = 0
6x − 3dx = 0, a solution may exist. We find u
x y =  12
−1n − 1 1 cos
3x cosh
3
 − y + cos
nx cosh
ny + C. 3 −3 sinh
3 n=1  n sinh
n   2 2  n 
−1n + 6
1 −
−1n  sin
ny cosh
n
1 − x u
x y = 2 n4  4 sinh
 n=1 1 R  r 2 u
r  = a0 +
2 cos2  − 1 2 2 R 1
ln
y2 +
 − x2 e− sin
d + c u
x y = 2 −
2
a cos
xe−y d + c, with a = − f
 cos
d u
x y =  0 0

1. Since

0

CHAPTER 20 Section 20.1 1. 26 − 18i

3.

1
1 + 18i 65 n 4n+1

11. i4n =

i2 2 n = 1 = 1 i 13. 23. 27. 29.

5. 4 + 228i

7. 6 − i

1
−1632 + 2024i 4225 4n 2 4n+3 4n 2

9.

= ii4n = i, since i4n = 1 i4n+2 = i
i  = −1 i = i i i = −i    + 2n 19.  − tan−1 23 + 2n 17. a2 − b2 + b + 1 2ab − a 21.  + 2n 2 √  √      25. 29 cos − tan−1 25 + i sin − tan−1 25 2 2 cos
3/4 + i sin
3/4 √  −1  1     + i sin tan−1 18 65 cos tan 8          z−w   z−w  =  = 1  z − w  = 1. Hint: If z = 1, then z¯z = 1 and  1 − z¯ w   z¯z − w¯z  ¯z  z − w 

Section 20.2 1. 5. 9. 13.

15. 17.

19.

21.

√ Circle of radius 9 with center (8, −4) 3. Circle of radius 21 65 with center (0, − 21 ) The real axis for x ≤ 0 7. The line y = x + 2 The line 8x + 10y + 27 = 0 11. The half-plane 3x + y + 2 > 0 K is the closed half-plane 2x + 8y + 15 ≥ 0; every point of K is a limit point of K, and there are no other limit points; boundary points of K are those points on the line 2x + 8y + 15 = 0; K is closed but not compact (because K is not bounded). M consists of all points below the line y = 7; limit points are points of M and points on the line y = 7; boundary points are points x + 7i; M is open; M is not compact. U consists of all points x + iy with 1 < x ≤ 3; limit points are points of U and points on the line x = 1; boundary points are points on the lines x = 1 and x = 3; U is neither open nor closed; U is not compact (neither closed nor bounded). W consists of all x = iy with x > y2 . These are points (x y) inside and to the right of the parabola x = y2 . This set is open, and not compact. Limit points are all points of W and points on the parabola; boundary points are the points on the parabola. 1 + 2i 23. 2 − i 25. −1 27. 23 i

29. If n is even, say n = 2m, then i2n = i4m = 1, so {1} is one convergent subsequence; if n = 2m + 1, then i2n = i4m i2 = i2 = −1, so {−1} is another convergent subsequence. There are others.

Answers and Solutions to Selected Problems

A62

CHAPTER 21 Section 21.1 1. u = x v = y − 1; Cauchy–Riemann equations hold at all points; f is differentiable for all complex z.  3. u = x2 + y2  v = 0; nowhere; nowhere 5. u = 0 v = x2 + y2 ; the Cauchy–Riemann equations hold at (0, 0); nowhere 7. u = 1 y/x; nowhere; nowhere 9. u = x2 − y2  v = −2xy;
0 0; nowhere y x  v = −4y − 2 ; Cauchy–Riemann equations hold for all nonzero z; differentiable for z = 0 11. u = −4x + 2 x + y2 x + y2

Section 21.2 1. 2; z + 3i < 2 3. 1; z − 1 + 3i < 1 5. 2; z + 8i < 2 7. 1; z + 6 + 2i < 1 9. No; i is closer to 2i than 0 is, so if the series converged at 0 it would have to converge at i. 11. cn+1 /cn is either 2 or 21 , depending on whether n is odd or even, so cn+1 /cn has no limit. However, lim cn 1/n = 1, so the radius of convergence is 1 by the nth root test applied to



n→

cn z . n

n=0

Section 21.3 1. cos
1 + i sin
1 7.

3. cos
3 cosh
2 − i sin
3 sinh
2

1

1 − cos
2 cosh
2 + 21 i sin
2 sinh
2 2

13. u =

9. i

sin
x cos
x

5. e5 cos
2 + ie5 sin
2 2

2

11. u = ex −y cos
2xy v = ex cosh
y sinh
y

2 −y 2

sin
2xy

v = cos2 cosh2
y + sin2
x sinh2
y cos2
x cosh2
y + sin2
x sinh2
y 2  2  1 iz 1 iz 15. sin2
z + cos2
z =
e − e−iz +
e + e−iz = 1 2i 2  + 2n  n any integer 19. z = ln
2 +
2n + 1i n any integer 17. z = ln
2 + i 2

Section 21.4 3 4n − 1 i ln
4 + i 3. ln
5 +
2n + 1i ln
5 + i 1. ln
4 + 2 2 √ √

     5. ln
85 +
2n + 1 + tan−1 − 29 i ln
85 +  + tan−1 − 29 i   z  z  i
arg
z−arg
w =  e  7. Hint: In polar form, w w

Section 21.5 1. ie

−
2n+/2

−
2n+/2

3
2n+3/4

    3ln
2 3ln
2 − i sin cos 2 2

3. e 5. e   n n 7. cos + + i sin + 9. 16e
2n+1 cos
ln
4 − i sin
ln
4  8 2  8 2
2n + 1
2n + 1 + i sin 13. cos
n/3 + i sin
n/3 11. 2 cos 4 4 

15. The nth roots of unity are k = e2ki/n for k = 0 1     n − 1. Now use the fact that for z = 1 z = e2i/n .

n−1  k=0

zk =

zn − 1 , with z−1

Answers and Solutions to Selected Problems CHAPTER 22 Section 22.1 The graph of the curves in Problems 1,3,5,7, and 9 are shown below. 1. initial point 6 − 2i, terminal point 2 − 2i; simple and not closed; tangent $

t = 2ieit = −2 sin
ti + 2 cos
tj

y 0

3. 1 + i 3 + 9i; simple and not closed; $

t = 1 + 2ti = i + 2tj

y

3

2

4

5

x

6


0.5
1.0
1.5
2.0 5. 3, 3; closed but not simple; )

t = −3 sin
t + 5 cos
ti = −3 sin
ti + 5 cos
tj

9 8 7 6 5 4 3 2 1 1.0

1.5

2.0


2

4 2
1

0
2

1

2

3

x


4

0
2
4
6
8
10
12
14
16

1 2

3

9. 1, cos
2 − 2 sin
4i; simple and not closed; $

t = − sin
t − 4 cos
2ti = − sin
ti − 4 cos
2tj

y 1.5 1.0 0.5
0.4
0.2 0
0.5

x

y
1


2

3.0

7. −2 − 4i 4 − 16i; simple and not closed; +

t = 1 − ti = i − tj

y


3

2.5

0.2 0.4 0.6 0.8

1.0

x

4

x

A63

Answers and Solutions to Selected Problems

A64

Section 22.2 1. 8 − 2i

3
1 + i 2

3.

1
−13 + 4i 2

5.

7. − 21 cosh
8 − cosh
2

9. − 21 e−1
cos
2 + i sin
2 − e
cos
2 − i sin
2 larger number)

11. 10 + 210i

13.

25 i 2

15.

2
1 + i 3

√ 16. 1/ z (or any

Section 22.3 1. 0

3. 0

5. 2i

7. 0

9. 0

11. 4i

Section 22.4 1. 32i 3. 2i
−8 + 7i 9. −512
1 − 2i cos
256

5. −2e2
cos
1 − i sin
1 − 39i 13. 2 11. − 13 2

7. i
6 cos
12 − 36 sin
12

CHAPTER 23 Section 23.1 
−1n 2n 2n 2 z  z < (that is, the series converges for all complex z) n=0
2n! √   1
z − 4in  z − 4i < 17 5.
n + 1zn  z < 1 3. n+1
1 − 4i n=0 n=0

1.

7. −3 +
1 − 2i
z − 2 + i +
z − 2 + i2  z < 9. 63 − 16i +
−16 + 2i
z − 1 − i +
z − 1 − i2  z < 
−1n
z + i2n+1  z < 11. n=0
2n + 1!   n  2n 2 + 2n−1 2n z +i z2n+1  z < 13. 1 + iz +
2n!
2n + 1! n=1 15. Fix z and think of w as the variable. Define f
w = ezw . Then f
n
w = zn ezw . By Cauchy’s integral formula, 1
ezw n!
ezw zn = dw, with $ the unit circle about the origin. Then dw, so f
n
0 = zn = 2i $ wn+1 n! 2i $ wn+1   n 2  n 2


 n n  z  1 z n ezw 1 1  1 z z z zw zw dw = e dw. Then = e dw = n! 2i $ n!wn+1 2i n=0 $ n!wn+1 2i $ n=0 n! w w n=0 n!
1 1 = ez
w+1/w dw. Now let w = ei on $ to derive the result. 2i $ w 17. The maximum must occur at a boundary point of the rectangle. Consider each side. On the left vertical side, x = 0 and ez  = eiy  = 1. On the right vertical side, ez  = e1 eiy  = e has maximum e. On the lower horizontal side, ez  = ex for 0 ≤ x ≤ 1, with maximum e. On the upper horizontal side, ez  = ex has maximum e. Thus the maximum of ez  on this rectangle is e.

Section 23.2 1.

 1
−1n +
z − in , for 0 < z − i < 2 z − i n=0
2in+1

5. −

1 − 2 −
z − 1 0 < z − 1 < z−1


−1n+1 4n 2n−2 z  z <
2n! n=1 1  2i z2n 7. 2 +  0 < z < 9. 1 +  0 < z − i < z z−i n=0
n + 1!

3.

CHAPTER 24 Section 24.1 1. pole of order 2 at z = 0

3. essential singularity at z = 0

7. simple pole at −i, removable singularity at i

5. simple poles at i and −i, pole of order 2 at 1
2n + 1 9. simple poles at 1, −1 i, and −i 11. simple poles at  2

Answers and Solutions to Selected Problems 13. f has a Taylor expansion f
z =



A65

an
z − z0 n in some open disk about z0 , and g has a Laurent expansion of the

n=0

 b b
z − z0 n in some annulus 0 < z − z0  < r. Then fg has an expansion of the form form g
z = −1 + z − z0 n=0 n b−1 a0  + c
z − z0 n in this annulus, and b−1 a0 = 0 because b−1 = 0 and a0 = f
z0  = 0. z − z0 n=0 n

1 g
z g
z = in some annulus h
z
z − z0 3 q
z  0 < z − z0  < r. Now g/q is analytic at z0 , and so has Taylor expansion g
z/q
z = cn
z − z0 n in some disk

15. Write h
z =
z − z0 3 q
z, where q is analytic at z0 and q
z0  = 0. Then

n=0

about z0 . Further, c0 = 0 because q
z0  = 0 and g
z0  = 0. Then, in some annulus about z0   c0 g
z = + c
z − z0 n−2 . h
z
z − z0 3 n=1 n

Section 24.2 and at −2i 251
9 + 12i; the value of the integral is therefore 2i.  8i
e − 1 3. 0 5. 2i 7. 2i 9. −i/4 11. 0 13. 2i 15. 2   an
z − z0 n and h
z = bn
z − z0 n , with a0 = 0 and b3 = 0. From Problem 15, Section 24.1, 18. Write g
z = n=0 n=3        g
z n n n n = d
z − z0  , with d−3 = 0. Write g
z = an
z − z0  = bn
z − z0  dn
z − z0  and h
z n=−3 n n=0 n=3 n=−3 equate the coefficient of
z − z0 n on the left with the coefficient of
z − z0 n in the product on the right. We get a0 = d−3 b3  a1 = d−3 b4 + d−2 b3  a2 = d−3 b5 + d−2 b4 + d−1 b3 . Use these to solve for d−1 in terms of coefficients 1 1 a0  a1  a2  b1      b4 and use the fact that an = g
n
z0  bn = h
n
z0 . n! n! 1. The residue at 1 is

1
16 − 12i, 25

Section 24.3 1. 2i 3. 0 1 −4t 5. cos
3t 7. 36 e − 361 e2t + 16 te2t 9. 21 t2 e−5t √ √ √ 
1 + 5e−4  11. 2/ 3 13. 2/3 15. 41 e−2 2 sin
2 2 17. −/128 19. 32   3 17 cos
/5 + 16 cos
3/5 16 sin 21. − 5 5 289 + 168 cos
2/5 + 136 cos
4/5 + 32 cos
6/5  
cos
x 1 − 23. Re s
eiz /
z2 + 1 i = − 21 ie− , so dx = 2i − ie = e− . 2 2 − x + 1
z dz, with $ the unit 25. With the trigonometric substitutions, we obtain −4i 2 2 2 2 2 2 2 $
 −  z + 2
 +  z +
 −   0 circle. The two poles within the unit disk are z = ±

− , and the residue at each is −i/2. Therefore, the value +

of the integral is 2i
−i/, or 2/.  −z2 e dz = 0, where $ is the rectangular path. Writing the integral over each piece of the $

R 2 2 e−x dx + e−
R+it idt+ boundary (starting on the bottom and going counterclockwise), we have −R 0
−R
0
 2 2 2 2 e−
x+i dx + e−
−R+it idt = 0. Let R → . The second integral is e−R et e−2iRt dt, and this goes to R  0
2 e−x dx− zero as R → . Similarly, the fourth integral has limit zero. The first and third integrals give −


√ 2 2 2 2 2 2 e−x e e−2ix dx = 0. Then  = e e−x cos
2xdx + ie− e−x sin
2xdx. The last integral is zero − − − √ 2  2 because the integrand is odd. Finally, write  = 2e 0 ex cos
2xdx, since this integrand is even. Now solve for the integral.

27. By Cauchy’s theorem,

Answers and Solutions to Selected Problems

A66




2

eiz dz = 0. Now integrate over each piece of $, going counterclockwise and beginning
/4 2 2i
0
R 2 i/4 2 eix dx + eiR e Riei d + ei
re  ei/4 dr = 0. The second integral tends to with the segment 0 R 

28. By Cauchy’s theorem,

$

0

0

R

zero as R → , and in the limit √ the last
equation becomes √
1 2 2 2 2 −r 2
1 + i
1 + i . Equate the real part of each side and the

cos
x  + i sin
x dx = e dr = 2 2 2 0 0 imaginary part of each side to evaluate Fresnel’s integrals.
2

1 1 1 z 30. dz = −4i d = dz. The integrand has double    2 2 2 + 2z + 2
 +  cos
 iz
z $ $ 0 1  z+ + 2 z   1 1 2 2 poles at
− ±  −  , but only
− + 2 − 2  is enclosed by the unit circle $. Compute   
2   i 1 −4iz 1 2 2 − 2 − + Re s =− 2   . Then d = 2 . 2 2 2 3/2
z + 2z +  
 −  
 +  cos
2
 − 2 3/2 0

2  1 1 d = 2 d. Finally, check that 2 2
 +  cos
 0 0
 +  cos


CHAPTER 25 Section 25.1 1. The images are given by the following diagrams.

v

y π D

w=ez

C

A

π

B

x (a)

C' –eπ

D' –1

A'

y 2

C'

C

D' –1 – π2

B

D

B'

C B 1

A

C'

θ=

π

D' A'

x (c)

B'

1

e

u

e A'

e2 B'

u

v

y π

u

(b) v

4

e

A'

y π

1 e

x

1

A

u

v π

D

B' eπ

1

D

C

A 1

B 2

x (d)

C'

D'

Answers and Solutions to Selected Problems v

y 2

D

e2

C'

π

C D'

–1

1 e

x

2

u

A' A

B

– π2

B'

(e)

3. The images are given by the following diagrams.

v

y π

w = 4 sin(z)

D

C

D'

2

A

π

(4, 0)

B

x

2

A'

(a)

y

B' v

π

D

C

A

B π

4

2

u

D'

u2 – v2 = 8

4

π

C'

(4, 0) x

A'

(b)

B'

C'

u

v

y D'

C' D

π

C

6

A

x

B

A'

(c) v

y π

D

C'

C

2

B

A π

3π 2

2

1

B' (–4, 0)

A' (4, 0)

C

A

B

u

(d) D'

v D

D'

x

y 2

u

B'

A' u

1

2

B'

x (e)

C'

A67

Answers and Solutions to Selected Problems

A68

5. The sector consisting of all w with an argument in /2 .     1 1 1 1 7. Hint: Put z = rei and obtain u = r+ cos
k v = r− sin
k, where the half-line is  = k. 2 r 2 r 9. The entire w plane with the origin excluded.

Section 25.2 16 − 16i +
−7 + 13iz 4 − 75i +
3 + 22iz 3 + 8i −
1 + 4iz 3. w = 5. w = 4 + 7i −
2 + 3iz 4 − 8i +
−1 + 2iz −21 + 4i +
2 + 3iz 11 1 208 19 377 v = 10 9.
u − 2 +
v + 2 = 11.
u − 12 +
v + 2 = 21 63 3969 4 16 w = z¯ reverses orientation. az + b dw − b and ad − bc = 0, then z = − is also a linear fractional transformation. If w = cz + d cw − a az + b has either one or two solutions, depending on whether or not c is zero. A translation has Hint: Show that z = cz + d no fixed point.

1. w = 7. 13. 15. 17.

In each of Problems 19 and 21, there are many solutions, of which one is given here. 1 1 1 19. z → → −4 → −4 + i = w 21. z → iz → iz −
2 − 7i = w z z z

Section 25.3 In Problems 1, 3 and 5, there are many mappings having the stated property. We give one such mapping in each case. −4i
z + 1 3z + 2 + 6i 5. w = 7. w = z1/3 1. w = 2z + 1 − i 3. w = z + 2i z−1
1 $
m$
n 9. Hint: Evaluate f
1 f
−1 f
0, and f
 and then use the result that for positive tm−1
1 − tn−1 dt = $
m + n 0 integers m and n.

Section 25.4

g
t y
dt, where u
x 0 = g
x  −
t − x2 + y2
2 1 g
x0 + R cos
t y0 + R sin
t 3. u
x y = 2 0

R2 −
x − x0 2 −
y − y0 2 × 2 dt R +
x − x0 2 +
y − y0 2 − 2R
x − x0  cos
t − 2R
y − y0  sin
t 1
2 r cos
t − r sin
t
1 − r 2  5. u
r cos
 r sin
 = dt 2 0 1 + r 2 − 2r cos
t − 
1 1
1 − t cos
t/2 7. u
x y = 8 −1 1 + sin2
t/2  4 sinh
x/2 cos
y/2 1 + sin2
t/2 + sinh
x sin
y 1 − sin
t/2 dt × sinh2
x/2 + sin2
y/2 − 2 cosh
x/2 sin
y/2 sin
t/2 + sin2
t/2 1. u
x y =

Section 25.5 1. With a = Kei , equipotential curves are 
x y = K x cos
 − y sin
 = constant. These are lines of the form y = cot
x + b. Streamlines are 
x y = K y cos
 + x sin
 = constant, which are lines y = − tan
x + b. Velocity = f

z = Ke−i . There are no sources or sinks. 3. 
x y = cos
x cosh
y 
x y = − sin
x sinh
y. Equipotential curves are graphs of y = cosh−1
K/ cos
x, streamlines are graphs of y = sinh−1
C/ sin
x. 5. 
x y = K lnz − z
0  
x y = K arg
z − z0 . Equipotential curves are circles z − z0  = r and streamlines are rays emanating from z0 

C

−vdx + udv = 2K, with C the circle of radius r about z0 .

Answers and Solutions to Selected Problems

A69

  x y x . Equipotential curves are graphs of x + 2 + i y − = c, streamlines are graphs x2 + y 2 x2 + y 2 x + y2 y of y − 2 = d. x + y2

b x arg
z = c. Streamlines are graphs of − 9. Equipotential curves are graphs of K x + 2 2 x + y 2

  y b ib 1
k y− 2 + + lnz = d. Stagnation points occur where f = 0, or
z = k 1 − x + y2 4 z2 2z ib b2 ± 1− . z=− 4k 16 2 k2

7. f
z = k x +

CHAPTER 26 Section 26.1 1. (a) abcdefghijkl, where each of the letters a,   , g can be either an H (for heads) or T (for tails), and each of the symbols h,   , l can be any of the integers from 1 through 6 inclusive. (b) 27 65 3. 610
15 5. There are 43 53 = 8000 outcomes, more than enough. 7. 7(15)(8)(17), or 14,280 9. 7(11)(6)(14), or 6,468

Section 26.2 1. 9! 3. 9! 5. (a) 7! (b) 6! (c) 5! 7. 6! 9. 12! 11. 6!

Section 26.3 1. 5. 7.

25 P7 ,

or 2,422,728,000 3. 22 P6 , or 53,721,360 or 5.7408(1016 ) if order counts; 52 C10 , or 158,200,024,220 if order is unimportant 9. with order: 52 P5 , or 311,875,200; without order, 52 C5 , or 2,598,960 20 C4 , or 4,845 52 P10 ,

Section 26.4 1. Typical outcome: a b c d, where each letter can be any of the integers 1, 2, 3, 4, 5, 6. There are 64 outcomes. A has 40 outcomes in it, and B has 16 outcomes. 3. Typical outcome:   ! , with   ! and  four distinct letters (English) alphabet. There are 26 C4 outcomes. A has 24 C2 outcomes, and W has 7 C4 outcomes. 5. Typical outcome: a b c d e f , where each of a b and c can be H or T , and each of d e and f can be any 1, 2, 3, 4, 5, 6. There are 23 63 outcomes. C has 10 outcomes, D has 80 outcomes, E has 24 outcomes. 7. Typical outcome: a b c d e, in which each letter can be any of 1, 2, 3, 4, 5, 6. U has 21 outcomes, K has 21 outcomes. 9. Typical outcome: a b     p, with each letter a distinct card; there are 52 C16 outcomes. Y has one outcome, M has 4(40 C5 ) outcomes.

Section 26.5 1. (a) 5/16 (b) 13/16 3. (a) 1/221 (b) 105/221 5. (a) 5/14 (b) 15/28 (c) 9/28 7. (a) 1/1860480 (b) 1/4 (c) 1271/1292 9. (a) 1892/23205 (b) 2257/54145

Section 26.6 1. 1 −

756 + 57 /67 , about 0.3302

3. 1 −8 C5 /52 C5 , about 0.9998

5. 1 − 1/341 055
365
364 · · ·
365 − n + 1 7. For n people, the probability of at least two having the same birthday is 1 − , assuming a 365n 365-day year.

A70

Answers and Solutions to Selected Problems

Section 26.7 1. (a) 1/2 (b) 2/3

3. (a) 5/16 (b) 5/11

5. (a) 56/1296 (b) 4/671

7. 1/249,900

9. 1/16, 1/4, 0

Section 26.8 1. 5. 9. 11.

dependent 3. dependent independent 7. dependent at least two heads, 0.4864, exactly two heads, 0.3456 probability of drawing exactly two reds is 108/113

Section 26.9 1. Pr(10) = 1/2, Pr(5) = Pr(50) = Pr(0) = Pr(1200) = 1/8 3. Pr(Ford) = 1/9, Pr(Chevy) = 1/18, Pr(Porsche) = 1/6, Pr(VW) = 1/12, Pr(Lamborghini) = 1/12, Pr(tricycle) = 1/12, Pr(Honda) = 1/18, Pr(Mercedes) = 1/9, Pr(tank) = 1/12, Pr(Steamer) = 1/24, Pr(bike) = 1/8 5. Pr(0) = 1/4, Pr(50, 000) = 1/20, Pr(.05) = Pr(1000) = 1/20 Pr(500) = Pr(1500) = 1/5, Pr(20) = Pr(lion) = 1/10

Section 26.10 1. Pr(≥10 yearsitem defective) = 0.14, Pr
5−10 yearsitem defective = 035, Pr
1−5 yearsitem defective = 012, Pr